Theory of Magnetism
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Section One: Reading Comprehension
Theory of Magnetism
To understand the magnetic behavior of materials, it is necessary to take a
microscopic view of matter. A suitable starting point is the composition of the
atom, which Bohr described as consisting of a heavy nucleus and a number of
electrons moving around the nucleus in specific orbits. Closer investigation
reveals that the atom of any substance experiences a torque when placed in a
magnetic field; this is called a magnetic moment. The resultant magnetic
moment of an atom depends upon three factors-the positive charge of the
nucleus spinning on its axis, the negative charge of the electron spinning on
its axis, and the effect of the electrons moving in their orbits. The magnetic
moment of the spin and orbital motions of the electron far exceeds that of the
spinning proton. However, this magnetic moment can be affected by the
presence of an adjacent atom. Accordingly, if two hydrogen atoms are
combined to form a hydrogen molecule, it is found that the electron spins, the
proton spins, and the orbital motions of the electrons of each atom oppose
each other so that a resultant magnetic moment of zero should be expected.
Although this is almost the case, experiment reveals that the relative permeability of hydrogen is not equal to 1 but rather is very slightly less than unity.
In other words, the molecular reaction is such that when hydrogen is the
medium there is a slight decrease in the magnetic field compared with free
space. This behavior occurs because there is a precessional motion of all
rotating charges about the field direction, and the effect of this precession is
to set up a field opposed to the applied field regardless of the direction of
spin or orbital motion. Materials in which this behavior manifests itself are
called diamagnetic for obvious reasons. Besides hydrogen, other materials
possessing this characteristic are silver and copper.
Continuing further with the hydrogen molecule, let us assume next that
it is made to lose an electron, thus yielding the hydrogen ion. Clearly,
complete neutralization of the spin and orbital electron motions no longer
takes place. In fact, when a magnetic field is applied, the ion is so oriented
that its net magnetic moment aligns itself with the field, thereby causing a
slight increase in flux density.This behavior is described as paramagnetism and
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is characteristic of such materials as aluminum and platinum. Paramagnetic materials
have a relative permeability slightly in excess of unity.
So far we have considered those elements whose magnetic properties
differ only
very slightly from those of free space. As a matter of fact the vast majority of materials
fall within this category. However, there is one class of materials-principally iron and its
alloys with nickel, cobalt, and alumi- num-for which the relative permeability is very
many times greater than that of free space. These materials are called ferromagnetic
and are of great importance in electrical engineering. We may ask at this point why iron
its alloys) is so very much more magnetic than other elements. Essentially, the
answer is provided by the domain theory of magnetism. Like all metals, iron is crystalline in
structure with the atoms arranged in a space lattice. However, domains are subcrystalline
particles of varying sizes and shapes containing about 10
atoms in a volume of
approximately cubic centimeters. The distinguishing feature of the domain is that the
magnetic moments of its constituent atoms are all aligned in the same direction Thus in a
ferromagnetic material, not only must there exist a magnetic moment due to a nonneutralized spin of an electron in an inner orbit, but also the resultant spin of all neighboring
atoms in the domain must be parallel.
Figure 1-1. Representation of a Ferromagnetic Crystal: (a) Unmagnetized and (b)
Fully Magnetized by the Field H
It would seem by the explanation so far that, if iron is composed of completely
magnetized domains, then the iron should be in a state of
complete magnetization
throughout the body of material even without the application of a magnetizing force.
Actually, this is not the case, because the domains act independently of each other, and for
a specimen of unmagnetized iron these domains are aligned haphazardly in all
directions so that the net magnetic moment is zero over the specimen. Figure 1-1 illustrates
situation diagrammatically in a simplified fashion. Because of the crystal
lattice structure of iron the ‘easy’ direction of domain alignment can take place in any
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one of six directions-left, right, up, down, out, or in-depending upon the direction of the
applied magnetizing force. Figure l-l(a) shows the unmagnetized configuration. Figure
l-l(b) depicts the result of applying a force from left to right of such magnitude as to
effect alignment of all the domains. When this state is reached the iron is said to be
saturated-there is no further increase in flux density over that of free space for further
increases in magnetizing force.
Large increases in the temperature of a magnetized piece of iron bring about a
decrease in its magnetizing capability. The temperature increase enforces the agitation
existing between atoms until at a temperature of 750°C the agitation is so severe that it
destroys the parallelism existing between
the magnetic moments of the
neighboring atoms of the domain and thereby causes it to lose its magnetic property. The
temperature at which this occurs is called the curie point.
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
…….1. With his atomic theory, Bohr contributed to the understanding of
the magnetic behavior of materials.
…….2. The atoms of a substance, if placed in a magnetic field, are subject
to a torque.
…….3. Platinum is a diamagnetic material.
…….4. In ferromagnetic materials, the magnetic moments of large groups
…….5. In an unagnetized ferromagnetic material, the domains are
aligned in different direction.
…….6. The magnetic properties of iron increase with an increase in
B. Choose a, b, c, or d which best completes each item.
1. Permeability of silver is less than unity .......... .
a. because of its atoms setting up a field against the applied field
b. because of its molecules rotating about the applied field
c. due to the precessional spin of its positive charges
d. due to the orbital motions of its negative charges
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2. It is true that .......... .
a. paramagnetic materials provide a small penetration of the magnetic
b. paramagnetic materials provide a great penetration of the magnetic
c. the resultant magnetic moment of an atom depends on its spinning
d. the resultant magnetic moment o f an atom depends on the nucleus
spinning on its axis
3. According to the text, .......... .
a. two atoms of hydrogen, if combined, pronounce a
greater than 1
b. two atoms of hydrogen, if combined, give rise to a high magnetic
c. diamagnetic materials have magnetic properties more than those of
free space
d. diamagnetic materials have magnetic properties less than those of
free space
4. Paramagnetism is based on the fact that the magnetic moment of a
paramagnetic material, when placed in a magnetic field, .......... .
a. results in a decrease in flux density
b. lines up with the field
c. is equal to 1
d. is low compared with free space
5. The magnetic properties of diamagnetic and paramagnetic materials
.......... those of free space.
a. are greater than
b. are smaller than
c. differ slightly from
d. differ greatly from
6. The abnormal magnetic properties of iron may be caused by ,.......... .
a. the magnetic moment resulting from an inner orbital spin of a
nonneutralized electron
b. the parallelism of the resultant spin of all neighboring atoms in the
c. the domains oriented at random with their axes pointing in various
d. both a and b
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C. Answer the following questions orally.
1. What is called a magnetic moment?
2. What does the resultant magnetic moment of an atom depend on?
3. How do adjacent atoms affect the magnetic moment of each other?
4. How does the magnetic behavior of materials differ?
5. Why does platinum have the characteristic of paramagnetism?
6. What forms the domains in a ferromagnetic material?
7. What causes the alignment of the magnetic domains in iron?
8. What is called the curie point?
Part IL Language Practice
A. Choose a, b, c, or d which best completes each item.
1. Copper is .......... material, therefore, it exhibits a relative permeability
slightly less than unity.
a. a paramagnetic
b. a diamagnetic
c. a permeable
d. a neutral
2. Iron provides a great penetration of the magnetic field, that is ,its
.......... is many times greater than that of free space.
a. magnetic flux
b. atomic composition
c. relative permeability
d. magnetic moment
3. Elements and metals which have slight magnetic properties are called
.......... materials.
a. magnetic
b. metallic
c. diamagnetic
d. paramagnetic
4. Iron and some of its alloys have an appreciable magnetic permeability.
These materials are called .......... .
a. ferromagnetic
b. diamagnetic
c. paramagnetic
d. magnetic
5. The state of .......... is reached when all the magnetic domains are
aligned in one direction.
a. magnetization
b. saturation
c. flux density
d. neutralization
B. Fill in the blanks with the
1. Magnet
form of the words
a. Maxwell showed that some of the properties of ..........
compared to a flow.
may be
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b. Lines of flux are conventionally said to leave a .......... material at the
north pole and re-enter at the south pole,
c. If the .......... field is produced by a solenoid, we will have the same
representation of lines of flux, but with the solenoid taking the
place of a ......... .
2. Permeate
a. Relative a pure number that is the same in all unit systems;
the value and dimension of absolute ..........depend upon the system
of units employed.
b. A .......... is an apparatus used for determining corresponding values
of magnetizing force and flux density in a test specimen.
3. Move
a. When a conductor is .......... through a magnetic field in such a way as
to cut the magnetic lines, an emf is generated in the conductor.
b. A moving - conductor microphone is a microphone the electric output
of which results from the .......... of a conductor in a magnetic field.
c. In a moving - conductor loudspeaker, the .......... conductor is in the
form of a coil connected to the source of electric energy.
4. Rotate
a. The most important parts of a dc motor are the ..........., the stator,
and the brushgear .
b. A .......…converter combines both motor and generator action in one
armature winding connected to both a commutator and slip rings,
and is exited by one magnetic field.
c. A rotary generator is an alternating-current generator adapted to be
.......... by a motor or prime mover.
5. Saturate
a. A magnetic-core
reactor operating in the region of saturation
without independent control means is known as .......... reactor.
b. A .......... sleeve is a flexible tubular product made from cotton and
coated with an electrical insulating material.
c. Saturation induction is sometimes referred to as .......... flux density.
C. Fill in the blanks with the following words.
Inductance is a characteristic of magnetic fields, and it was
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Faraday in his renowned experiments of 1831. In a general .......... inductance
can be characterized as that property of a circuit ......... by which energy is
capable of being stored in a magnetic .......... field. A significant and
distinguishing feature of inductance, however, is that .......... makes itself fell in
a circuit only when there is a/an ..........current or flux. Thus, although a circuit
element may have virtue of its geometrical and magnetic properties,
its presence in the .......... is not exhibited unless there is a time rate of change
of .......... .This aspect of inductance is particularly stressed when we consider
it .......... the circuit viewpoint. However, for the sake of completeness,
inductance is also.......... from an energy and a physical view.
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. Trans formers are to be found in such varied applications as radio and
television receivers and electrical power distribution circuits.
b. An understanding of electromagnetism is essential to the study of
electrical engineering because it is the key to the operation of a great
part of the electrical apparatus found in industry as well as the home.
c. Similarly, static transformers provide the means for converting energy
from one electrical system to another through the medium of a
magnetic field.
d. Other important devices-for example, circuit breakers, automatic switches, relays, and magnetic amplifiers-require the presence of a confined magnetic field for their proper operation.
e. All electric motors and generators, ranging in size from the fractional
horsepower units found in home appliances to the 25,000-hp giants
used in some industries, depend upon the electromagnetic field as the
coupling device permitting interchange of energy between an electrical
system and a mechanical system and vice versa.
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Section Two: Further Reading
The main experimental facts underlying magnetism are the following:
The ancient Greeks knew that the mineral loadstone or magnetite
(Fe3O4) attracts
pieces of iron at some zones on its surface. Magnetite in fact is a natural magnet.
If a piece of magnetite is brought near a bar of hard iron, this too acquires the
property of attracting iron filings; it has become an artificial magnet. This process is
known as magnetic induction.
When iron bars become magnetized, the quality of attracting pieces of iron is found
at two regions at the ends of the bars. These regions are called
the poles of the magnet.
If we bring two magnets together, with one magnet fixed and the other free to turn,
we see that the first magnet exerts some forces on the second. A magnet produces a
magnetic field in the space around it. In a similar way, we have seen that electric charges
produce an electric field.
The fact that a magnetic bar or a compass needle comes to rest in a roughly
north-south direction when freely suspended near the surface of the earth is evidence that
the earth itself acts as a magnet. By convention we give the name north pole to that pole
of the magnetic bar or needle which seeks the geographic north, the other pole being
known as the south pole.
If we take a given pole of a magnet and place it first at one and then at
the other
pole of a second magnet, in one case the two poles will attract each other and in the other
case they will repel each other. It is found that unlike poles attract whereas like poles
repel each other.
In a magnet it is not possible to separate the north pole from the south pole. In fact,
if we break a magnet in two we find a south pole at the broken end of the part that had
the original north pole, and a north pole at the broken end of the part that had the
original south pole.
Of all the metals or elements, only iron, cobalt and nickel and some of their alloys
have pronounced magnetic properties. These materials are known
as ferromagnetic
materials. Other elements and metals have slight magnetic properties, and they are called
paramagnetic materials. There is a third series
of materials that have magnetic
properties less than those of a vacuum, and these are called diamagnetic materials.
Ousted in 1820 showed that a current flowing in an electric circuit
exerts forces on
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a nearby magnet and so demonstrated that a magnetic field is generated around an electric
current. Consequently, if we place an electric circuit in a magnetic field, the circuit is
subject to forces.
The fact that a magnetic field can be produced either by a magnet or by an electric
current may seem strange. But we must remember that in matter we have microscopic
circuits due to the movement of electrons, and these circuits are responsible for the
magnetic effects of ferromagnetic materials. However, the causes which underlie the
magnetic forces produced by electric circuits are not fully understood (just as there are
still problems in our understanding of the forces between electric charges and the
nature of the force of gravity), although we know the laws that govern their actions and
can therefore use them. We know that atoms consist of a heavy central positive nucleus
and a number of electrons, in either circular or elliptical orbits, around the nucleus.
Recently there has been added the concept that each electron itself is spinning about an
axis through its centre, this motion being known as electron spin. Here, it is impossible to
offer a complete explanation
of this and we must limit ourselves lo saying that the
fundamental magnetic particles in ferromagnetic materials are the spinning electrons.
these elec - trons occupy definite shells in the atom, and some spin in one direction and
some in the other. Their magnetic effects tend to neutralize each other partially but not
wholly. The excess of those spinning in one direction over those spinning in the other
causes each atom as a whole to act as a small permanent magnet. Moreover, in
ferromagnetic materials there is the
existence of some kinds of interatomic forces that
cause the alignment of all magnetic effects of large groups of atoms to give highly
magnetic domains. In
an unmagnetized ferromagnetic substance these domains are
oriented at random with their magnetic axes pointing in various directions, so that the
resultant magnetic effect is zero. The application of an external field lines up the domain
axes, thereby giving rise to the magnetic effect of a ferromagnetic material.
In hard iron the domains do not easily return to their previous positions when the
external field is removed, while in soft iron this occurs fairly readily. Paramagnetic and
diamagnetic materials, on the other hand, are substances in which the arrangement of the
spinning electrons does not give appreciable magnetic properties. When the
temperature of a ferromagnetic material is raised beyond a certain value (known as
the Curie point), thermal agitation
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destroys the alignment within the domains and the materials lose their
ferromagnetic properties. These properties return when the materials are
cooled. The Curie point for iron is of the order of 700°C. As in the case of an
electric field, a magnetic field at each point may be defined by its field
strength. This is represented by the vector H. The direction is that in which a
north pole subjected to this field tends to move. Because the magnetic field
may be produced by a current, the strength can be defined in terms of current.
In order to do this we consider a solenoid, i.e., a coil of wire wound uniformly
on a cylindrical former. If the solenoid is long compared with its radius, we
can consider that a uniform magnetic field is produced inside the coil, parallel
to its axis. If N is the number of turns, I the length of the solenoid and I the
current that flows in the coil, we have H=NI/I .The magnitude of H is
measured in amperes per metre, and the quantity NI is expressed in amperes.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1.If we break a magnetic bar into two pieces, the two poles at the point
of breakage will ……. .
a. be two north poles
b. be a north pole and a south pole
c. pronounce greater attraction
d. pronounce smaller attraction
2.It is true that ……. .
a. the circuits in matter produce magnetic forces
b. the circuits in an electric current produce magnetic forces
c. paramagnetic materials have smaller magnetic properties than diamagnetic materials
d. diamagnetic materials have greater magnetic properties than ferromagnetic materials
3. The factors bringing about the magnetic properties of materials are the
spinning ……. .
a. nuclei
c. neutrons
b. atoms
d. electrons
4. Paragraph ten mainly discusses ……. .
a. the magnetic field
b. the electric field
c. the theory of magnetism
d. the theory of gravity
5. When the temperature of cobalt is below the Curie point ….... .
a. all magnetism disappears
b. some magnetism disappears
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c. the metal has appreciable magnetic properties
d. the alignment of the magnetic domains is destroyed
6. The vector H representing the field strength of a magnetic field may be
expressed as the product of ……. .
a. the number of turns in a coil and the current in amperes which
flows through it
b. the number of turns in a coil and the current in amperes which flows
through it per unit length
c. the current flowing through a coil and the length of the coil
d. the current flowing through a coil per unit length
7. In a bar magnet, the magnetic domains …….. .
a. neutralize each other
b. repel each other
c. are at random
d. are aligned
8. A magnet and an electric current in a circuit produce a magnetic field
by virtue of …….. .
a. the position of the magnetic domains
b. the orientation of the atomic nuclei
c. the movement of the electrons
d. the alignment of the interatomic forces
B. Write the answers to the following questions.
1. How have the poles of a magnetic bar been initially named?
2. How do you describe the process of magnetic induction?
3. How is an electric field compared with a magnetic field?
4. How may a magnetic field be demonstrated?
5. How are paramagnetic materials different from diamagneti materials?
6. What did Oersted prove in 1820?
7. What is the difference between the soft iron and the ilard iron?
8. How do you describe the vector H?
Section Three: Translation Activities
A. Translate the following passage into Persian.
A mongst the oldest and easiest to observe scientific phenomena are those of
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magnetism. The subject of bar magnets and magnetic poles is given the name
of magnetostatics by analogy with electrostatics. Magnetic phenomena,such as
poles and the fields they produce, can be explained in terms of fields due to
electric currents. Thinking on the atomic scale, the electrons which circulate
round the heavy central positive nucleus constitute a current. So each orbiting
electron produces a magnetic field. In general, the orbits of the electrons are
disposed in random planes in space and so the net magnetic field is zero.
Should a suitable stimulus be applied, the orbits can be aligned so that their
magnetic fields are in the same direction. In some materials the orbits, once
aligned, stay that way and these are the materials which produce permanent
magnets. In other materials the orbits return to their random dispositions
once the stimulus is removed-these are the materials used as electromagnets.
For a given stimulus, they produce a greater field than the materials used for
permanent magnets. 1t should be mentioned at this stage that ability to
produce a high field is not the only factor to be considered when deciding
upon a material to be used for an electromagnet; there are problems of energy
loss to be considered if the state of magnetization is to be changed frequently.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
2. alloy
3. apparatus
5. brush gear
6. circuit
7. commutator
8. compass needle
9. compose
10. coupling device
11. crystalline
12. depict
13. diamagnetic
14. dispose
15. domain
16. electrostatic
17. exert on
18. ferromagnetic
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19. gravity
20. lattice
21. magnetic amplifier
22. magnetic induction
23. magnetic moment
24. manifest
25. mineral
26. neutralize
27. nucleus
28. orbit
29. permeability
30. precession
32. repel
33. saturate
34. sleeve
35. specimen
37. stator
38. stimulus
41. transformer
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Section One: Reading Comprehension
Power Stations
are five sources of energy which together account for nearly all the
world's electricity. They are coal, oil, natural gas, hydroelectric power and
nuclear energy. Coal, oil and nuclear plants use the steam cycle to turn heat into
electrical energy, in the following way. The steam power station uses very pure
water in a closed cycle. First it is heated in the boilers to produce steam
at high
pressure and high temperature, typically 150 atmospheres and 550°C
in a
modern station. This high-pressure steam drives the turbines which in turn
drive the electric generators, to which they are directly coupled. The maximum
amount of energy will be transferred from the steam to the turbines only if the
latter are allowed to exhaust at a very low pressure, ideally a vacuum. This can
be achieved by condensing the outlet steam into water. The water is then pumped
back into the boilers and the cycle begins again. At the condensing stage a large
quantity of heat has to be extracted from the system. This heat is removed in the
condenser which is a form of heat exchanger. A much larger quantity of cold
impure water enters one side of the condenser
and leaves as warm water,
having extracted enough heat from the exhaust steam to condense it back into
water. At no point must the two water systems mix. At a coastal site the warmed
impure water is simply returned to the sea
at a point a short distance away. A
2 GW station needs about 60 tons of sea water each second. This is no problem
on the coast, but inland very few sites could supply so much water all the year
round. The alternative is to
recirculate the impure water. Cooling towers are
used to cool the impure water so that it can be returned to the condensers, the
same water being cycled continuously. A cooling tower is the familiar concrete
structure like a very broad chimney and acts in a similar way, in that it induces
a natural draught. A large volume of air is drawn in round the base and leaves
through the open top. The warm, impure water is sprayed into the interior of
the tower from a large number of fine jets, and as it falls it is cooled by the rising
air, finally being collected in a pond under the tower. The cooling tower is really
a second heat exchanger where the heat in the impure water is passed
to the
atmospheric air; but unlike the first heat exchanger, the two fluids are
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allowed to come into contact and as a consequence some of the water is lost by
The cooling towers are never able to reduce the impure water temperature right down to the ambient air temperature, so that the efficiency of the
condenser and hence the efficiency of the whole station is reduced slightly
compared with a coastal site. The construction of the cooling towers also
increases the capital cost of building the power station. The need for cooling
water is an important factor in the choice of sites for coal, oil and nuclear
plants. A site which is suitable for a power station using one type of fuel is
not necessarily suitable for a station using another fuel.
Coal-Fired Power Stations
Early coal-burning stations were built near the load they supplied. A station of
2 GW output, consumes about 5 million tons of coal in a year. In Britain where
most power station coal is carried by rail, this represents an average of about 13
trains a day each carrying 1000 tons. This means that large coal-fired stations need
a rail link unless the station is built right at the pit head.
Oil-Fired Power Stations
Power station oil can be divided into crude oil which is oil as it comes from the
well, and residual oil which remains when the more valuable fractions have
been extracted in the oil refinery. The cost of moving oil by pipeline is less than
that of moving coal by rail, but even so stations burning crude oil are often
sited near deep-water berths suitable for unloading medium-sized tankers.
Stations burning residual oil need to be sited near to the refinery which supplies
them. This is because residual oil is very viscous and can only be moved through
pipelines economically if it is kept warm.
Nuclear Power Stations
In contrast to coal and oil the cost of transporting nuclear fuel is negligible
Because of the very small amount used . A 1 GW station needs about 4 1 2 tons
of uranium each week. This compares very favourably with the 50,000 tons of
fuel which would be burnt each week in a comparable coal-fired power
station. Present nuclear stations use rather more cooling water than comparable coal-fired or oil-fired plants due to their lower efficiency. All nuclear
stations in Britain, with one exception, are situated on the coast and use sea
water for cooling.
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Hydroelectric Power Stations
Hydroelectric power stations must be sited where the head of water is available,
and as this is often in mountainous areas, they may need long transmission lines to
carry the power to the nearest load center or link up
with the grid. All
hydroelectric schemes depend on two fundamental factors: a
flow of water and a
difference in level or head. The necessary head may be obtained between a lake and a
nearby valley, or by building a small dam in a river which diverts the flow through
the power station, or by building a high
dam across a valley to create an artificial
Part I. Comprehension Exercises
A. Put "T" for true and "F" for false statements . Justify your
........ 1. Gas and nuclear plants use the steam cycle to turn heat
........ 2. Condensers remove the heat from the outlet steam.
........ 3. The steam power station uses pure water in an open cycle.
........ 4. Steam pressure affects the generators directly.
........ 5. Having cooled off the exhaust steam, the warmed impure water
may be recirculated.
........ 6. Natural air is forced through the cooling tower.
........ 7. Large coal-fired stations situated far from the pit head need a rail
........ 8. Oil-fired power stations consume certain constituents of crude oil.
........ 9. Nuclear power stations use less cooling water than comparable
coal-fired or oil-fired plants.
........ 10. Hydroelectric power stations have to be built where there is
enough water pressure.
B. Choose a, b, c, or d which best completes each item.
1. In steam power stations, the turbine efficiency will increase
a. the steam pressure is kept constant
b. the outlet steam is condensed into water
c. the steam temperature is not varied
d. the outlet water is pumped back into the boilers
2. The steam power station uses pure water .......... .
a. to produce the steam required to drive the turbines
b. to produce the steam required to activate the generators
if .......... .
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c. to create the vacuum space necessary for the system
d. to create the pressure and temperature needed
3. The heat of the steam is removed by .......... the condenser.
a. the recirculation of cold pure water in
b. the flow of natural air in one side of
c. the recirculation of the steam in
d. the flow of cold water through one side of
4. Prior to recirculation, impure water must be cooled ........... .
a. in broad concrete structures
b. in broad metal chimneys
c. at the bottom of the tower
d. at the top of the tower
5. The cooling factor in a cooling tower is .......... the tower.
a. the pond under
b. the interior of
c. the water inside
d. the air passing through
6. Systems recirculating impure water, compared with those on the coast,
……… .
a. decrease the efficiency of the station
b. increase the capital cost of building the station
c. reduce the impure water temperature to the required level
d. both a and b
7. The first paragraph mainly discusses .......... .
a. the structure of a condenser compared with that of a cooling
b. the mechanism of the steam power station
c. the main sources of energy which account for electricity
d. the cooling water as a deciding factor in the choice of sites for coal,
oil, and nuclear plants
C. Answer the following questions orally.
1. What are the five sources of energy used for the generation of electrical
2. What are the two water systems used in the condenser?
3. What is the water resulted from steam condensation used for?
4. How much sea water does a 2 GW station need each second?
5. How is the mechanism of a cooling tower similar to that of a
6. How do you describe the mechanism of a cooling tower?
7. What are the two heat exchangers used in the system?
8. How much coal does a 2 GW station consume every year?
9. Why should stations burning residual oil be sited near to the refinery
which supplies them?
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10. Why is the cost of transporting nuclear fuel negligible compared with
coal and oil?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. The energy of water may be converted to work by hydraulic .......... .
a. turbines
b. generator
c. boilers
d. towers
2. Gas oil must be .......... and then used.
a. isolated
b. heated
c. refined
d. vapourized
3. In the condenser, the outlet steam is .......... and recirculated.
a. exchanged
b. condensed
c. depressurized
d. purified
4. Cooling towers cause water to be .......... .
a. condensed
b. exhausted
c. evaporated
d. recycled
5. Air pump suction must be applied to the lowest pressure point or
points within a condenser which are normally at the inlet tube plate
where .......... rate and hence steam side pressure drop are greatest.
a. the condensation
b. the temperature
c. the cooling
d. the evaporation
B. Fill in the blanks with the appropriate form of the words
1. Exchange
a. The atomic movements of materials are said to be held in parallel or
antiparallel by exchange forces, thought to be due to the sharing or
.......... of electrons between neighbouring atoms in the crystal
structure of the material.
b. Coupling forces, similar to the .......... forces of the atom, exist
between the molecules of a compound.
c. Cooling towers and condensers are two kinds of heat .......... .
2. Circulate
a. A .......... register retains data by inserting it into a delaying means
and regenerating and reinserting the data into the register.
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b. A constant flow of electrolyte through a cell to facilitate the
maintenance of uniform conditions of electrolysis is known as ..........
of electrolyte.
C. A................magnetic wave is a traverse magnetic wave for which the
lines of magnetic force form concentric circles.
3. Couple
a. Water heated in the boilers of the steam power station produces
steam at high pressure which drives the turbines .......... to generators.
b. Typical oscillators are in practice amplifiers in which power is fed
into the grid circuit from the plate circuit by means of either
electrostatic or electromagnetic .......... between these circuits.
4. Condense
a. Condensed-mercury temperature is the temperature measured on the
outside of the tube envelope in the region where the mercury is .........
in a glass tube or at a designated point on a metal tube.
b. Steam can be .......... into water.
c. A .......... is a form of heat exchanger.
5. Drive
a. A .......... is an electronic circuit that supplies input to another
electronic circuit.
b. Grid .......... power is the average of the product of the instantaneous
values of the alternating components of the grid current and the
grid voltage over a complete cycle .
b. A system consisting of one or several electric motors and of the entire
electric control equipment designed to govern the performance of these motors is called the electric .......... .
C. Fill in the blanks with the following words.
Any steam power station burning coal or .......... could fairly easily be
converted to burn .......... gas. Such stations must, of course, be .......... near a
large gas main. However, it is .......... felt that natural gas is too high a/an ..........
fuel and too valuable as an industrial ..........stock and home heating fuel, to
be.......... in power stations. The point is that gas burnt to produce electricity
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which might then be used for home heating, would produce .......... at about 33
percent efficiency, whereas the .......... gas burnt in a domestic boiler would
.......... heat at up to 80 percent efficiency.
D. Put the following sentences in the right order to form a paragraph.
Write the corresponding letters in the boxes provided.
a. The flame temperature is clearly much higher than the steam temperature, but the thermodynamic efficiency of a conventional station
depends on the steam temperature not the flame temperature.
b. Firstly, work associated with existing coal- and oil-burning power
stations where efforts are being made to utilize the inherent thermodynamic efficiency of the very high flame temperatures of burning oil or
pulverised coal.
c. Generators have been constructed to convert some of the energy in the
flame, which is a moving ionised gas, directly into electricity.
d. Experiments and design studies are being carried out to develop new
ways of generating electricity.
e. Secondly, work is being done to try to convert solar energy into
f. These fall broadly into three groups.
g. Thirdly, we have what is sometimes called the nuclear alternative,
h. These are known as magnetohydrodynamic generators.
Section Two: Further Reading
Electric Power Stations
Characteristics Influencing Generation and Transmission
There are four main characteristics of electricity supply which, however obvious,
have a profound effect on the manner in which it is engineered. They
are as
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(a) Electricity, unlike gas and water, cannot be stored and the supplier
has little control over the load at any time. The control engineers endeavor to keep the
output from the generators equal to the connected load at the
specified voltage and frequency.
(b) There is a continuous increase in the demand for power. Although in
industrialized countries the rate of increase has declined in recent years,
even the modest rate entails massive additions to the existing systems. A large
and continuous process of adding to the system thus exists. Networks are
evolved over the years rather than planned in a clear-cut manner and then left
(c) The distribution and nature of the fuel available. This aspect is of great interest
as coal is mined in areas not necessarily the main load centres; hydroelectric power is
usually remote from the large load centres. The problem of station siting and
transmission distances is an involved exercise in economics. The greater use of nuclear
energy will tend to modify the existing pattern of supply.
(d) In recent years environmental considerations have assumed major importance
and influence the siting, construction cost, and operation of generating plants. Planning
is also affected because of delays in making a start to projects because of legal proceedings,
etc. Of particular importance at the present time is the question of the environmental
impact of nuclear plants, especially the proposed fast breeder reactor.
Energy Conversion Employing Steam
The combustion of coal or oil in boilers produces steam at high temperatures and
pressures which is passed to steam turbines. Oil has economic advantages when it can be
pumped from the refinery through pipelines direct to the boilers of the generating
station. The use of energy resulting from nuclear fission is being progressively extended
in electricity generation; here also the basic energy is used to produce steam for turbines.
The axial-flow type of turbine is in common use with several cylinders on the same shaft.
The steam power station operates on the Rankine cycle, modified to
include muted superheating, feed-waterheating, and steam reheating. Increased
thermal efficiency results from the use of steam at the highest possible
pressure and temperature .Also , for turbines to be economically constructed
500 MW and over a re now being used. With steam turbines of 100 MW
capacity and over the efficiency is increased by reheating the steam after it has
been partially expanded , by an external heater . The reheated steam is then
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Figure returned to the turbine where it is expanded through the final stages of blading.
A schematic diagram of a coal-fired station is shown in Figure 2-1. In
Figure 2-2, the
flow of energy in a modern steam station is shown. Despite
continual advances
in the design of boilers and in the development of
improved materials, the nature
of the steam cycle is such that efficiencies are comparatively low and vast quantities of
heat are lost in the condensate.
However, the great advances in design and materials in
the last few years have
increased the thermal efficiencies of coal stations to about 40
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In coal-fired stations, coal is conveyed to a mill and crushed into fine
powder i.e., pulverised. The pulverized fuel is blown into the boiler where it
mixes with a supply of air for combustion. The exhaust from the LP turbine is
cooled to form condensate by the passage through the condenser of large
quantities of sea- or river-water. Where this is not possible cooling towers are
Fluidized-Bed Boilers. For typical coals, combustion gases contain 0.2-0.3
percent sulphur dioxide by volume. If the gas flow-rate through the granular
bed of a great-type boiler is increased the gravity pull is balanced by the
upward gas force and the fuel-bed lakes on the character of a fluid. In a
travelling grate this increases the heat output and temperature. The ash
formed conglomerates and sinks into the grate and is carried to the ash pit.
The bed is limited to the ash-sintering temperature of 1050-1200°C. Secondary
combustion occurs above the bed where CO burns to CO2 and H2S to SO2.
type of boiler is undergoing extensive development and is attractive because
of the lower pollutant level and better efficiency.
Energy Conversion Using Water
Perhaps the oldest form of energy conversion is by the use of water power. In
the hydroelectric station the energy is obtained free of cost. This attractive
feature has always been somewhat offset by the very high capital cost of
construction, especially of the civil engineering works. Today, however, the
capital cost per kilowatt of hydroelectric stations is becoming comparable with
that of steam stations. Unfortunately, the geographical conditions necessary for
hydro-generation are not commonly found. In most highly developed countries
hydroelectric resources are used to the utmost.
An alternative to the conventional use of water energy, pumped storage,
enables water to be used in situations which would not be amenable to
conventional schemes. The utilization of the energy in tidal flows in channels
has long been the subject of speculation. The technical and economic difficulties
are very great and few locations exist where such a scheme would be feasible.
An installation using tidal flow has been constructed on the La Rance estuary in
northern France where the tidal height range is 9.2 m (30 ft) and the tidal flow is
estimated at 18,000 m3/s.
Before discussing the types of turbine used, a brief comment on the
general modes of operation of hydroelectric stations will be given. The vertical
difference between the upper reservoir and the level of the turbines is known
head. The water falling through this head gains kinetic energy which it
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then imparts to the turbine blades. There are three main types of installation as
(a) High. Head or Stored-the storage area or reservoir normally fills in over
400 h;
(b) Medium Head or Pondage-storage Tills in 200-400 h;
(c) Run of River-storage fills in less than 2 h and has 3-15 m head.
A schematic diagram for type (c) is shown in Figure 2-3.
Associated with these various heights or heads of water level above the turbines
are particular types of turbine. These are:
(a) Pelton. This is used for heads of 184-1840 m (600-6000 it) and consists
of a bucket wheel rotor with adjustable How nozzles.
(b) Francis. Used for heads of 37-490 m (120-1600 ft) and is of the mixed flow
(c) Kaplan. Used for run of river and pondage stations with heads of up to 61 m
(200 ft). This type has an axial-flow rotor with variable-pitch blades.
Figure 2-3. Hydroelectric Scheme —Kainji, Nigeria. Section through the intake dam
and power house. The scheme comprises an initial four 80 MW Kaplan turbine sets
with the later installation of eight more sets. Running speed 115.4 rev/min. This is a
large-flow scheme with penstocks 9 m in diameter. (Permission of Engineering.)
Typical efficiency curves for each type of turbine are shown in Figure
2-4. As the efficiency depends upon the head of water which is continually
fluctuating, often water consumption in cubic meters per kilowatt-hour is used and
is related to the head of water. Hydroelectric plant has the ability to start UP
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quickly and the advantage that no losses are incurred when at a standstill.
has great advantages, therefore, for generation to meet peak loads at minimum
cost, working in conjunction with thermal station. By using remote control of the
hydro sets, the time from the instruction to start up to the actual connexion to
the power network can be as short as 2 min
Figure 2-4. Typical Efficiency Curves of Hydraulic Turbines.
Gas Turbines
The use of the gas turbine as a prime mover has certain advantages over steam
plant, although with normal running it is less economical to operate. The main
advantage lies in the ability to start and take up load quickly. Hence the gas turbine is
coming into use as a method for dealing with the peaks of the system load. A
further use for this type of machine is as a synchronous compensator to assist with
maintaining voltage levels. Even on economic grounds it is probably advantageous
to meet peak loads by starting up gas turbines from cold in the order of 2 min rather
than running spare steam plant continuously.
Comprehension Exercises
Choose a, b, c, or d which best completes each item.
1. We may deduce from the text that .......... .
a. since electricity cannot be stored, enough electricity must be
generated at all times to meet the variations in demand
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b. since the supplier does not have control over the load, variations in
demand have to be limited to certain degrees
c. gas, water, and electricity can be stored to satisfy the unexpected
increases in demand
d. gas, water, and power systems are not amenable to ordinary energy
2. It is true that .......... .
a. coal is usually mined in accessible areas
b. in many industrialized countries the rate of increase in power
demand has declined in recent years
c. environmental considerations do not have any effect on the siting
and operation of power plants
d. industrialized countries do not add any new networks to their
systems due to a decline in power demand
3. In order to increase the efficiency of a steam power station, .......... .
a. steam turbines of 100 MW capacity must be employed
b. coal and oil must be used at high temperatures and pressures
c. the reheated steam must be expanded and returned to the turbine
d. the steam must be used at the highest possible pressure and
4. It is true that .......... .
a. the advances in the design of boilers have not affected the efficiency
of coal stations
b. coal stations have low efficiencies because of the heat lost in the
steam cycle
c. systems operating on the steam cycle have high efficiencies
d. the MW capacity of all the steam turbines used today is over 500
5. In fluidized-bed boilers …..... .
a. the upward gas force causes the fuel-bed to take the character of a
b. the fluid characteristic of the fuel-bed increases the heat output
c. CO burns to CO2 and H2 S to SO2.
d. all of the above
6. The use of energy in tidal flows .......... .
a. has greatly replaced the conventional use of water energy
b. has enabled man to make use of water energy wherever he likes
c. may be an alternative to the conventional use of water energy
d. is a common way of using water energy without any difficulties
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7. The last paragraph mainly discusses .......... .
a. how the gas turbine deals with the peaks of the system load
b. how the gas turbine is used as a synchronous machine
c. the advantages of the gas turbine
d. the mechanism of the gas turbine
B Write the answers to the following questions.
1 What are the four main characteristics of electricity supply?
2. What are coal and nuclear energy used for?
3. What is the function of an external heater?
4. What process does the coal go through in coal-fired stations?
5. When are cooling towers used in coal-fired systems?
6. Why are fluidized-bed boilers called so?
7. What are the advantages of fluidized-bed boilers?
8. What is the most prominent feature of hydroelectric power stations?
9. What are the initial requirements for hydro-generation?
10. How does water obtain the energy required to impart to the turbine
11. What are the three types of installation?
12. What is the head of water?
13. How is Kaplan turbine different from Francis and Pelton turbines?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Magnetohydrodynamic (MHD) Generation
In conventional power generation, fuel such as oil or coal is burned. The burning
fuel heats boilers to produce steam. The steam is used to drive turbo-alternators.
The MHD process generates electricity without requiring a boiler or a turbine.
MHD generation works on the principle that when a conductor cuts a Magnetic
field, a current flows through the conductor. In MHD generation the conductor is an
ionized gas. Small amounts of metal are added to the gas to improve its
conductivity. This is called seeding the gas. The seeded gas is then
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pumped at a high temperature and pressure through a strong magnetic field. The
electrons in the gas are collected at an electrode. This movement of
electrons constitutes a current flow.
Two methods of MHD generation can be used: the open-cycle and the
closed-cycle. In the open-cycle method the hot gas is discharged. In the
closed-cycle method it is recirculated.
The open-cycle method uses gas from burning coal or oil. The gas is
and then passed through a magnetic field to generate current. The seeding elements are
recovered and the gas can then be used to drive a turbine before being allowed to
The closed-cycle method uses an inert gas, such as helium, which is heated
indirectly. The gas is circulated continually through the MHD
MHD generation is still in its early stages but already an efficiency rate of 60% has
been reached. This compares with a maximum of 40% from conventional power
B. Find the Persian equivalents of the following terms and
Expressions and write them in the spaces provided.
1. ambient
2. breeder reactor
3. chimney
4. concrete
5. condense
6. conglomerate
7. convey
8. crude oil
9. decline
11. efficiency
12. email
13. exchange
16. fluctuate
17. grate type
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18. incur
19. inland
20. magnetohydrodynamic generator
21. outlet
22. pit head
23. prominent
24. pulverize
25. seeded gas
26. speculate
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Section One:Reading Comprehension
Electrical Insulation
Insulation is required to keep electrical conductors separated from each other
and from
other nearby objects. Ideally, insulation should be totally nonconducting, for then
currents are totally restricted to the intended conductors. However, insulation does
conduct some current and so must be regarded as a material of very high resistivity. In
many applications, the current flow due to conduction through the insulation is so small
that it may be entirely neglected. In some instances the conduction currents, measured
by very sensitive instruments, serve as a test to determine the suitability of the
insulation for use in service.
Although insulating materials are very stable under ordinary circum stances, they
may change radically in characteristics under extreme conditions
of voltage stress or
temperature or under the action of certain chemicals. Such changes may, in local regions,
result in the insulating material becoming highly conductive. Unwanted current flow
brings about intense heating and the rapid destruction of the insulating material. These
insulation failures account for a high percentage of the equipment troubles on
electric-power systems. The selection of proper materials, the choice of proper shapes
and dimensions, and the control of destructive agencies are some of the problems of the
insulation-system designer.
Many different materials are used as insulation on electric-power systems. The
choice of material is dictated by the requirements of the parti- cular application and by
cost. In residences, the conductors used in branch circuits and in the cords to appliances
may be insulated with rubber or plastics
of several different kinds. Such materials can
withstand necessary bending, are relatively stable in characteristics, and are inexpensive.
They are subjected to relatively low electrical stress.
High-voltage cables are subjected to extreme voltage stress; in some cases several
hundred kilovolts are impressed across a few centimeters of insulation. They must be
manufactured in long sections, and must be
sufficiently flexible as to permit pulling into
ducts of small cross section. The insulation may be oil-impregnated paper, varnished
cambric, or synthetic materials such as polyethylene.
The coils of generators and motors may be insulated with tapes of various kinds.
Some of these are made of thin sheets of mica held together by a binder, and other are
of fiber glass impregnated with insulating varnish.
This insulation must be capable of withstanding quite high operating tempera-mechanical
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forces, and vibration.
The insulation on power-transformer windings is commonly paper tape
and pressboard operated under oil. The oil saturates the paper, greatly
its insulation strength, and, by circulating through ducts, serves as
an agent for
carrying away the heat generated due to 12R losses and core losses in the transformer. The
transformer insulation is subjected to high electric
stress and to large mechanical
forces. The shape and arrangement of conducting metal parts is of particular concern
in transformer design.
Overhead lines are supported on porcelain insulators. Between the supports air
serves as insulation. Porcelain is chosen because of its resistance
to deterioration when
exposed to the weather, its high dielectric strength, and
its ability to wash clean in rain.
Part L Comprehension Exercises
A. Put "T" for true and "F" for false statements. Justify your answers.
........ 1. The higher the insulation the less the loss of power.
........ 2. In order to avoid insulation failures, very expensive materials are
used in power systems.
........ 3. Insulation failures do not affect the electric equipment.
........ 4. Voltage and temperature variations may bring about insulation
........ 5. Rubber and plastic insulating materials are preferred to other
kinds because of their cost.
......... 6. Polyethylene and mica have different applications in electricalpower systems.
B. Choose a, b, c, or d which best completes each item.
1. The first paragraph mainly discusses .......... .
a. electrical conductors
b. nonconducting materials
c. the purpose of insulation
d. the application of insulators
2. As we understand from the text, .......... .
a. perfect insulation is not possible
b. stable insulators are not available
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c. chemicals do not affect good insulators
d. insulators may never change to temporary conductors
3. The second paragraph mainly discusses .......... .
a. the problems caused by the insulation-system designer
b. the factors resulting in insulation failures
c. the characteristics of insulating materials
d. the rapid destruction of insulating materials
4. Tapes of insulating fiber glass are commonly used to insulate ......…..
a. ordinary conductors
b. the windings of power transformer
c. high-voltage cables
d. the coils of generators and motors
5. Insulating tapes .......... .
a. cannot withstand high electrical stress
b. can withstand high temperatures
c. are used to insulate ducts of small cross section
d. are used to stop deterioration caused by the weather
C. Answer the following questions orally.
1. What is the purpose of insulation?
2. What is one way of deciding on the suitability of insulators used for different
3. What does the choice of an insulating material depend on?
4. What is the application of porcelain insulators?
5. What are insulating tapes used for?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. In order to keep electrical conductors separated from each other, ……. .
materials must be used.
a. capacitive
b. resistive
c. insulating
d. conducting
2. The .......... of metals increases with increase of temperature.
a. conductivity
b. resistivity
c. solubility
d. durability
3. Voltage stress may affect .......... of insulating materials.
a. the sensitivity
b. the suitability
c. the stability
d. the conductivity
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4. Certain insulating materials are impregnated with oil; that is, they are
……… oil.
a. saturated with
b. covered with
c. deprived of
d. made of
5. Porcelain has high resistance to deterioration; in other words, it does
not......... quickly.
a .deflect
d. decline
C. decrease
B. Fill in the blanks with the appropriate form of the words given.
1. Arrange
a.Modern types of air blast or oil circuit breakers have frequently been
fitted with trip-free mechanisms in which the tripping instructions are
.......... to override the closing instructions.
b. Different insulation .......... have different characteristics.
c.On important lines of high lightning incidence, it is accepted practice
to seek to prevent direct strokes to phase conductors by .......... one
or more shielding wires above the phase conductors to intercept
lightning strokes and conduct them to the ground.
2. Operate
a. The design of a power system should be such that, when breakdowns
are inevitable, they are confined to locations where they cause
minimum damage and the least disturbance to .......... .
b. The whole of the electrical and mechanical quantities that characterize the work of a machine at a given time is known as ...........
c. A switch can be .......... by a lever or other operating means.
3. Subject
a. If the insulation were .......... to the normal operating voltage which
varies within quite narrow limits, there would be no problem.
b. It has been known for many years that an insulator surface .......... to
high voltage (dc or pulsed) in vacuum can acquire a large positive
c. The parameters which determine the lightning performance of a
transmission line are .......... to large variations according to frequency
distribution laws which can only be determined by field observations.
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4. Measure
a. The gas pressure can be ..……… by means of a standard pressure
b. Most branches of science and technology rely on electrical ........…. for
the control of processes and machines as well as for information.
c. The role played by electrical ........…..instruments is vital to all modern
laboratories and factories.
d. The .......... current is the value read on the microammeter during a
direct high-voltage test of insulation.
5. Insulate
a. For the economic transmission of power over considerable distances
the voltage
must be high, although with higher voltages the …........
cost rises.
b. When any object is said to be ........... , it is understood to be .......... in
suitable manner for the conditions to which it is subjected.
c. An insulated joint is used to .......... adjacent pieces of conduits, pipes,
rods, or bars.
d. The solid .......... generally are of the form of annular discs and
truncated cones.
C. Fill in the blanks with the following words.
due to
Modern porcelain insulator are designed and ……… in such a fashion
that in themselves they are almost .......... in operation. Flashover of line
insulators is almost always .......... to the breakdown of the air around them
……… overvoltage from lightning or other .......... . Insulators whose surfaces
are contaminated and then .......... by light rain or fog may flash over even
under ......…-operating-voltage conditions. If an insulator is ...…… or porous
and permits lighting or power- frequency, .......... to pass through the body of
the insulator, it may be .......... with the resultant dropping of the line.
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. Under voltag stress, an increase in temperature causes an increase of
the conductance of the insulating material.
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b. The increased temperature causes an increase in the conductance of
that local area, and more current flows, thereby increasing the temperature even more.
c This is sometimes called thermal breakdown.
d. Of course, the increase of temperature will cause increased conduction
of heat away from the region and a stable condition may result.
e. Conduction current flow implies a release of energy in the insulation.
f. However, if the voltage is further increased, the temperature may
continue to rise and the current may continue to increase until
chemical changes destroy the insulation and puncture results.
g. Under electrical stress near the puncture value, a local region may
increase in temperature because heat is released faster than it is carried
Section Two: Further Reading
Insulation Behavior
When insulation is placed between two metallic
conductors A and B connected to a voltage source
(Figure 3-1), several phenomena associated with
the insulation may be identified. The insulation, or
dielectric, influences the capacitance between the
plates, a current of low magnitude flows through
the body of the insulation,
a leakage current flows over the surface of the
insulation, and if the voltage is great enough,
sudden changes in the body of the insulation may
make it highly conductive.
Figure 3-1. Insulation Under Electrical
Stress Between Two Metal Electrodes
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Capacitance and Dielectric Hysteresis
As is well known, the presence of a dielectric between two conducting plates increases
the measured capacitance between these plates. The behavior of the electrons and
protons comprising the dielectric accounts for this capacitance increase. This
phenomenon is worth investigating, for it explains some other characteristics of
insulation of interest.
All matter is made up of protons, electrons, and neutrons. In the normal
state, these particles are grouped as atoms
or molecules in which the number of
electrons and the number of protons are
equal, therefore, each group is electrically
neutral. However, each of these particles
experiences a force due to its interaction
with any charges placed on nearby plates.
We say that the charged particles respond
to the electric field set up between the
In a perfect insulator, the electrons and
protons are held together in the atoms and
molecules and are not free to drift from one plate
to another. However, in the presence of an electric
field, they may move very slight distances, the
electrons toward the anode, the protons toward the
cathode. This situation is illustrated in greatly
simplified form in Figure 3-2a. This diagram
shows a group of polar molecules. Each of these is
neutral, but each has an extra electron at one end
and an extra proton at the other. These may be
moving with respect to each other but at some
instant have a position as shown. Let us confine
our attention to the polar molecule at P.
Next let charges be put on A and B by
connection to a voltage source as shown in Figure
3-2b. The molecule P rotates, taking up a new
position with
the electron displaced toward A and the proton toward 5, under influence of the electric
The effect of the many, many electron-proton pairs in the dielectric volume,
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moving as shown by the example P, is to produce an effect shown in Figure 3-2c.
Adjacent to the anode A there are an excess number of electrons
in the dielectric; near the cathode B are an excess number of protons in the
dielectric. These charges partially neutralize the effect of the charge originally placed on
the plate and additional charges move from the voltage source to A and B as the charges
in the dielectric take on their new positions. Hence, for the same voltage between the
plates, the charges that have moved from the source to the plates have been increased as
a result of the presence of the dielectric. The capacitance between the plates is greater,
therefore, than it would be without the dielectric.
Whenever a system of particles is moved from one position to another (such as
system P from the position shown in Figure 3-2a to that of Figure 3-2b), there are
forces which restrict the motion, and time is required to make the change. Such is the
case in dielectrics. In some materials the change is made in a fraction of a microsecond;
in others it may take several hours. During the period of change, the capacitance
appears to increase and current flows in the external circuit. It is sometimes stated that
charge is 'soaking' into the dielectric. The phenomenon is known as dielectric
When the voltage source is
disconnected from plates A and B and the
voltage between them is made zero by a
short-circuiting connection (Figure 3-3a),
the displaced particles in the dielectric
lend to go back to their normal state.
However, if it took a long time to get
them oriented, it will take a long time to
get them back into the normal state.
Hence the condition shown inFigure3-3a
may persist for some time, a charge
remaining on each plate equal in effect to
the charge remaining in the dielectric
adjacent to the plates
Figure 3-3. Charge Conditions in a Capacitor That
Next, suppose that the short
has Been Energized, (a) Immediately alter being
short-circuited, (b) Alter the short circuit
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circuit is removed. Forces continue to restore the dielectric to its neutral
state. With the circuit open, the charges held on the plates cannot be removed. As
the dielectric returns to its normal condition, the trapped charges on the plates produce a
voltage between A and B. This voltage may be serious hazard to a workman who expected
the capacitance between the plates to be discharged by the short-time application of a
short circuit. This
hazard is particularly serious on equipment of high
capacitance such as
high-voltage cables, static capacitors, and generator windings. For
this reason,
it is always desirable to keep such equipment continuously
short-circuited when workmen are to be in physical contact with the presumably
deenergized equipment.
Referring again to Figure 3-2a and b, the movement of particles, such as
the polar
molecule P, may result in the movement of other nonpolar molecules. If the molecular
motion is increased, the temperature of the material is increased. If the power supply is
an ac source, each reversal of voltage will tend to cause a reversal of the position of the
polar molecules and electrical energy from the source will be converted to heat in the
insulation. This loss is known as dielectric hysteresis. It increases with frequency and with
applied voltage. It must be considered in high-voltage cable design.
Conduction Currents
When voltage is applied between two plates separated by a dielectric (Figure 3-1), those
few free electrons that are present in the insulation drift from cathode to anode, This is
termed a conduction current (from anode to
cathode) and represents power loss into
the insulation. In insulation, the number of free electrons is low, and as a result the
resistivity of the material
is high, The number of free electrons may be increased by
an increase in temperature.
Surface Leakage Currents
Leakage currents flow along paths between electrodes over the surface of the insulating
material. The magnitude of these currents is in no way related to the resistivity of the
material itself. The value of the leakage current depends
on the applied voltage, the
insulation material, the surface contamination, and the moisture content of the air. On
seriously contaminated high-voltageline insulator surfaces leakage currents may be as
much as 100 milliamperes.
Insulation Breakdown
Insulation may undergo a very sudden change in characteristics in a process known as
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breakdown. Consider the arrangement shown in Figure 3-4. Two parallel-plane
electrodes A and B are separated by a sheet of dielectric of thickness t. A variable
voltage source V provides a difference of potential between A and B. Suppose the voltage
is slowly raised. At first the conduction current is very low, perhaps measurable in
microamperes. With increased applied voltage, the current suddenly increases, and the
insulation takes on the character of a metallic conductor. This is termed insulation
breakdown. On examination, a small damaged place may be found extending through the
insulation sheet. Perhaps there will be some charring and perhaps there will be a hole.
Figure 3-4. Insulation of Thickness t Being Stressed by Applied Voltage V.
The voltage at which such breakdown occurs is called the breakdown, or puncture,
voltage, VPs and the electric field intensity  p at that point is known as the breakdown
gradient or puncture strength of the insulation, where t is the insulation thickness.
The puncture strength of a particular sample is not a constant but varies
the thickness of the insulation, the shape and geometry of the electrodes, and the rate
of application of voltage.
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Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. The first paragraph mainly describes .......... .
a. the phenomenon of capacitance affected by a dielectric
b. the phenomenon of leakage current produced by a dielectric
c. the dielectric behavior when placed between two metallic conductors
connected to a voltage source
d. the dielectric connecting a voltage source to two metallic conductors
2. It is true that .......... .
a. the presence of an electric field causes the electrons and protons of a
dielectric to move
b. in a perfect insulator, the electrons and protons will never be
influenced by any external factor
c. the atomic particles of a dielectric will not interact with an electric
field placed close to it
d. in a dielectric, the electrons and protons are not very tightly held
together in the atoms
3. Paragraph five mainly describes .......... .
a. the behavior of electrons and protons comprising a dielectric
b. the characteristics of a dielectric placed between two conducting
c. why the electrons in the atoms of the dielectric move toward the
anode and the protons move toward the cathode
d. why the presence of a dielectric between two conducting plates
increases the capacitance between the plates
4. It is true that .......... .
a. the time required for systems of atomic particles to change positions
varies from one dielectric to another
b. some internal forces usually cause the systems of atomic particles to
move and change position
c. the capacitance between the plates does not change as the atomic
particles change positions
d. some insulating materials when disconnected from the voltage source
do not return to their normal state
5. Having removed the short circuit, equipments of high capacitance will
………. .
a. discharge the trapped charges on the plates very quickly and
cause the dielectric to return to its normal state
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b. be seriously dangerous to anybody in physical contact with them
c. displace the particles in the dielectric and cause a serious hazard to
the workmen
d. be unable to remove the charges on the plates and cause the
dielectric to go back to its normal state
6. As we understand from the text, .......... .
a. the number of free electrons in an insulating material decreases with
an increase in temperature
b. the resistivity of an insulating material decreases with an increase in
c. extreme temperature will permanently lower insulation resistance
d. dielectric hysteresis will cause loss of heat in the dielectric
7. It is true that .......... .
a. the lower the resistivity of a dielectric the higher the magnitude of
leakage currents
b. the higher the resistivity of a dielectric the lower the magnitude of
leakage currents
c. the value of the leakage current depends on factors such as surface
contamination and air moisture
d. the value of the leakage depends on factors such as temperature and
B. Write the answers to the following questions.
1. How does a dielectric affect the capacitance between two conducting
2. What may cause the electrons and protons of an insulator to move?
3. What is the phenomenon of dielectric absorption?
4. What happens if the voltage source is disconnected from the plates and
a short-circuiting connection is made?
5. What will happen if the short circuit is removed?
6. Why does the temperature of a dielectric between two conducting
plates increase?
7. What is dielectric hysteresis?
8. How does power loss into the dielectric occur?
9. How docs temperature affect an insulator as compared with a metal?
10. What is insulation breakdown?
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Section Three: Translation Activities
A. Translate the following passage into Persian.
Dielectric Heating
Dielectric heating is a method of heating a nonconducting material, a
dielectric, by high-frequency voltages. The material is placed between metal
plates across which a high-frequency supply is connected as shown in Figure
3-5. The dielectric and the plates then form a capacitor and an electrostatic
field is set up in the dielectric. As very high frequencies are used, up to 200
MHz, the movement of electrons in the dielectric becomes rapid. This causes
considerable heat in the substance.
Dielectric heating has two great
advantages over other forms of heating: it
provides rapid heat, and the heat is pro duced uniformly throughout the material.
Figure 3-5.
In other words, the inside of the material
gets hot at the same time as the surface. In addition, dielectric heating can be
controlled and it is predictable. Accurate heating times can be
calculated knowing the
dielectric properties of the materials to be heated.
Dielectric heating has many different uses, from the manufacture of
plastic raincoats
to baking biscuits. It is especially used in plastics, woodworking, and food industries.
A typical use is the manufacture of plywood. In the past, the layers of
wood and glue
were steam-heated under pressure until the glue melted and
the wood was firmly bonded.
The heat took a long time to penetrate the
wood, the glue did not melt uniformly and it
dried unevenly. With dielectric
heating, because of the difference in dielectric properties,
the glue melts
before the wood heats. It heats uniformly and it dries evenly. Using the
dielectric process, a single press can prepare 100 3-ply, 1 cm thick sheets of
plywood in about
30 minutes.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1. absorption
2. annular disc
3. breakdown gradient
4. circumstance
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5. contaminate
6. deprive
7. destruct g. deteriorate
9. dielectric
10. disturbance
11. duct
12. hysteresis
13. insulate
14. leakage
16. oil-impregnated paper 17.overvoltage
18. puncture
19. shatter
20. shielding wire
21. soak
22. solubility
23. stroke
24. trace
25. truncated cone
26. varnished cambric
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Section One: Reading Comprehens
The Distribution System
Although there is no ‘typical’ electric power system,
a diagram including the
several components that are usually to be found in the makeup of such a
system is shown in Figure 4-1; particular attention should be paid to those
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elements which will make up the component under discussion, the distribution system.
While the energy flow is obviously from the power generating plant to
the consumer, it may be more informative for our purposes to reverse the
direction of observation and consider events from the consumer back to the
generating source.
Energy is consumed by users at a nominal utilization voltage that may
range generally from 110 to 125 V, and from 220 to 250 V, the nominal
figures are 277 and 480 V, It flows through a metering device that determines
the billing for the consumer, but which may also serve to obtain data useful
later for planning, design, and operating purposes. The metering equipment
usually includes a means of disconnecting the consumer from the incoming
supply should this become necessary for any reason.
The energy flows through conductors to the meter from the secondary
mains (if any); these conductors are referred to as the consumer's service, or
sometimes also as the service drop.
Several services are connected to the secondary mains; the secondary
mains now serve as a path to the several services from the distribution
transformers which supply them.
At the transformer, the voltage of the energy being delivered is reduced
to the utilization voltage values from higher primary line voltages that may
range from 2200 V to as high as 46,000 V.
The transformer is protected from overloads and faults by fuses or
so-called weak links on the high-voltage side; the latter also usually include
circuit-breaking devices on the low-voltage side. These operate to disconnect
the transformer in the event of overloads or faults. The circuit breakers
(where they exist) on the secondary, or low-voltage, side operate only if the
condition is caused by faults or overloads in the secondary mains, services, or
consumers' premises; the primary fuse or weak link, in addition, operates in
the event of a failure within the transformer itself.
If the transformer is situated on an overhead system, it is also protected
from lightning or line voltage surges by a surge arrester, which drains the
voltage surge to ground before it can do damage to the transformer.
The transformer is connected to the primary circuit, which may be a
lateral or spur consisting of one phase of the usual three-phase primary main.
This is done usually through a line or sectionalizing fuse, whose function is to
disconnect the lateral from the main in the event of fault or overload in the
lateral. The lateral conductors carry the sum of the energy components
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flowing through each of the transformers, which represent not only the energy
used by the consumers connected thereto, but also the energy lost in the lines
and transformers to that point.
The three-phase main may consist of several three-phase branches connected
together, sometimes through other line or sectionalizing fuses, but sometimes also
through switches. Each of the branches may have several single-phase laterals
connected to it through line or sectionalizing fuses.
Where single-phase or three-phase overhead lines run for any considerable
distance without distribution transformer installations connected to them, surge
arresters may be installed on the lines for protection.
Some three-phase laterals may sometimes also be connected to the three-phase
main through circuit reclosers. The recloser acts to disconnect the lateral from the main
should a fault occur on the lateral, much as a line or sectionalizing fuse. However, it
acts to reconnect the lateral to the main, reenergizing it one or more times after a time
delay in a predetermined sequence before remaining open permanently. This is done so
that a fault which may be only of a temporary nature, such as a tree limb falling on the
line, will not cause a prolonged interruption of service to the consumers connected to
the lateral.
The three-phase mains emanate from a distribution substation, supplied from a
bus in that station. The three-phase mains, usually referred to as a circuit or feeder, are
connected to the bus through a protective circuit breaker and sometimes a voltage
regulator. The voltage regulator is usually a modified form of a transformer and serves
to maintain outgoing voltage within a predetermined band or range on the circuit or
feeder as its load varies. It is sometimes placed electrically in the substation circuit so
that it regulates the voltage of the entire bus rather than a single outgoing circuit or
feeder, and sometimes along the route of a feeder for partial feeder regulation. The
circuit breaker in the feeder acts to disconnect that feeder from the bus in the event of
overload or fault on the outgoing or distribution feeder.
The substation bus usually supplies several distribution feeders and carries the
sum of the energy supplied to each of the distribution feeders connected to it. In turn,
the bus is supplied through one or more transformers and associated circuit breaker
protection. These substation transformers step down the voltage of their supply circuit,
usually called the subtransmission system, which operates at voltages usually from
23,000 to 138,000 V.
The subtransmission systems may supply several distribution substations and may
act as tie feeders between two or more substations that are either of
the bulk power or transmission type or of the distribution type. They may also
be tapped to supply some distribution load, usually through a circuit breaker,
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for a single consumer, generally an industrial plant or a commercial consumer
having a substantially large load.
The transmission or bulk power substation serves much the same purposes as a distribution substation, except that, as the name implies, it handles
much greater amounts of energy: the sum of the energy individually supplied
to the subtransmission lines and associated distribution substations and losses.
Voltages at the transmission substations are reduced to outgoing subtransmission line voltages from transmission voltages that may range from 69,000 to
upwards of 750,000 V.
The transmission lines usually emanate from another substation associated with
a power generating plant. This last substation operates in much the same manner as
other substations, but serves to step up to transmission line voltage values the voltages
produced by the generators. Because of material and insulation limitations, generator
voltages may range from a few thousand volts for older and smaller units to some
20,000 volts for more recent, larger ones. Both buses and transformers in these
substations are protected by circuit breakers, surge arresters, and other protective
In all the systems described, conductors should be large enough that the energy
loss in them will not be excessive, nor the loss in voltage so great that normal nominal
voltage ranges at the consumers' services cannot be maintained.
In some instances, voltage regulators and capacitors are installed at strategic
points on overhead primary circuits as a means of compensating for voltage drops or
losses, and incidentally help in holding down energy losses in the conductors.
In many of the distribution system arrangements, some of the several elements
between the generating plant and the consumer may not be necessary. In a relatively
small area, such as a small town, that is served by a power plant situated in or very near
the service area, the distribution feeder
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may emanate directly from the power plant bus, and all other elements may
be eliminated, as indicated in Figure 4-2. This is perhaps one extreme; in
many other instances only some of the other elements may not be necessary;
e.g., a similar small area somewhat distant from the generating plant may find
it necessary to install a distribution substation supplied by a transmission line
of appropriate voltage only.
In the case of areas of high load density and rather severe service
requirements, the distribution system becomes more complex and
expensive. The several secondary mains to which the consumers’ services are connected
may all be connected into a mesh or network. The transformers supplying these
secondary mains or network are supplied from several
different primary
feeders, so that if one or more of these feeders is out of
service for any reason,
the secondary network is supplied from the remaining
ones and service to
the consumers is not interrupted. To prevent a
feeding-back from the
energized secondary network through the transformers connected to feeders out of
service (thereby energizing the primary and creating unsafe conditions),automatically
operated circuit breakers, called network protectors, are connected between the
secondary network and the
secondary of the transformers; these open when the
direction of energy flow is
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
........ 1. The text describes the distribution elements used between the
power generating plant and the consumer.
........ 2. The metering device is mainly used to offer data useful for design
and operating purposes.
........ 3. The three-phase main may consist of several three-phase branches
which in turn may consist of several single-phase laterals.
........ 4. Any power system must have secondary mains in order to supply
the consumer with energy through the services.
........ 5. Primary feeders are connected to the substation bus via voltage
regulators or protective circuit breakers.
........ 6. The voltage regulator is the same as the transformer.
........ 7. Subtransmission systems may be used as tie feeders between
different types of substations.
........ 8. The closer the power plant to the area it serves the fewer the
elements between the generating plant and the consumer.
9. Distribution transformers may directly supply the consumer with
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10. The elemental arrangement in a complex distribution system is so
that service interruption is improbable.
B. Choose a, b, c, or d which best completes each item.
1. It is true that .......... .
a. circuit breakers disconnect the high-voltage side of the transformer
in the event of overloads
b. fuses and circuit breakers are identical devices
c. fuses are weak links always installed on the secondary side of the
d. circuit breakers protect the high-voltage side of the transformer in
the event of overloads
2. The distribution transformer .......... .
a. is connected to the primary circuit through a line or sectionalizing
b. is connected to the primary main through a line or sectionalizing
c. helps to disconnect the lateral from the main in the event of fault or
d. helps to disconnect the consumers from the services in the event of
fault or overload
3. A circuit recloser is used to .......... .
a. connect a three-phase lateral to the three-phase main
b. disconnect the lateral from the main if a fault occurs on the lateral
c. reconnect the lateral to the main after a predetermined time delay
d. all of the above
4. It is true that .......... .
a. the substation bus supplies substation transformers
b. substation transformers supply the substation bus
c. the substation bus is supplied by the distribution feeders
d. substation transformers are the same as subtransmission systems
5. The bulk power substation handles the energy supplied to .......... .
a. the subtransmission lines and associated substations
b. the substation bus and the primary feeders
c. the distribution transformers and the meters
d. the secondary mains and associated services
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6. According to the text, .......... .
a. material and insulation limitations do not allow the generators to
work at their full capacities
b. material and insulation limitations have resulted in the use or
various protective devices to protect the generators
c. the voltages produced by the generators are stepped up in the
substation associated with the power plant
d. the voltages produced by the generators are reduced to usable
voltages in the substation associated with the power plant
7. The last paragraph mainly describes .......... .
a. consumers' service interruption
b. consumers’ service reliability
c. a complex distribution network d. a complex distribution system
C. Answer the following questions orally.
1. What is called the service?
2. What is the function of a surge arrester?
3. What does a lateral refer to?
4. What is the function of a voltage regulator?
5. What part does the substation bus play in the distribution system?
6. Where do the transmission lines originate from?
7. What is the function of a capacitor installed on an overhead primary
8. When does the distribution system become more complex?
9. What is a network?
10. How do network protectors help a distribution system?
Part IL Language Practice
A. Choose a, b, c, or d which best completes each item.
1. The .......... deliver electric energy from the secondary distribution or
street main, or other distribution feeder, or from the transformer, to
the wiring system of the premises served.
a. meters
b. buses
c. services
d. feeders
2. The function of .......... is to interrupt circuit faults.
a. a line
b. a service
c. a main
d. a transformer
3. A .......... serves as a common connection for two or more circuits.
a. fuse
b. switch
c. lateral
d. bus
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4. Two or more generating systems, substations, or feeding points may be
connected together by .......... .
a. a tie wire
b. a tie trunk
c. a tie feeder
d. a tie line
5. To automatically disconnect a transformer from a secondary network in
response to predetermined electric conditions on the primary feeder or
transformer, .......... are employed.
a. network relays
b. network protectors
c. circuit reclosers
d. circuit analyzers
B. Fill in the blanks with the appropriate form of the words given.
1. Connect
a. A connection diagram shows the .......... of an installation or its
component devices, controllers, and equipment.
b. A network Is .......... if there exists at least one path composed of
branches of the network, between every pair of nodes of the network.
c. A low voltage or secondary network is a continuous secondary main
or grid fed by a number of transformers ........... to the same primary
2. Protect
a. To ensure maximum ........... , the .......... system must possess a high
degree of electricity.
b. .......... equipment should be used against vibrations of voltage.
c. A differential relay responds to the difference between incoming and
outgoing electrical quantities associated with the ........... apparatus.
3. Limit
a. The function of a relay is to prevent or .......... damage during faults.
b. The inrush current of the rectifier transformer is generally the ..........
c. Hard limiting is a limiting action with negligible variation in output in the range
where the output is .......... .
d. A bridge ........... is a bridge circuit used as a limiter circuit.
4. Regulate
a. The substation may or may not require voltage .......... equipment.
b. The circuit on the output side of the .......... is known as the voltage
c. The voltage may be held constant at any selected point on the ..……..
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d. A voltage ......... relay is used on an automatically operated voltage
regulator to control the voltage of the regulated circuit.
5. Distribute
a. Electric power is received from substations and is .......... to the
consumers at voltage levels and degrees of continuity that are
acceptable to various types of consumers.
b. For a transverse electromagnetic wave on a two-conductor transmission line, the ......... constants are series resistance, series inductance
shunt conductance, and shunt capacitance per unit length of line.
c. A distribution switchboard is used for the distribution of electric
energy at voltages common for such .......... in a building.
d. A duct installed for occupancy of distribution mains is known
.......... duct
as a
C. Fill in the blanks with the following words.
Since a failure of a conductor results in a complete .......... to a circuit, it
is imperative that the causes of such .......... be minimized. The failure may
occur from .......... causes where the stresses and strains imposed are .......... too
great and the conductors literally tear .......... . More often, however, the cause
may initially be a/an .......... failure which then affects the conductor
mechanically. Overloads or .......... currents, for example, may cause heating of
the .......... to the point where they begin to liquefy and .......... mechanical
stresses can no longer be sustained and the conductors .......... apart, perhaps
vaporizing in the process.
D. Put the following sentences in the right order to form a
Write the corresponding letters in the boxes provided.
a. This accessibility and inaccessibility should prevail even under advert
or contingency conditions,
b. Moreover, the conductors and equipment on the poles should be so
situated that they can be handled safely by the people working on them
c. For example, a distribution pole line should be so located that free and
easy access to the facilities is available at all times, yet it should not
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interfere with pedestrian and vehicular traffic, nor intrude into areas
(such as playgrounds) where its presence may constitute a particular
d. Safe methods include the use of protective equipment, the use of
live-line tools and equipment, and the deenergization and grounding of
the facilities on which work is to be performed.
e. As much as practical, the utility's facilities have to be both accessible
(to the workers) and inaccessible (to the public).
f. This not only implies providing sufficient working space, but includes
considerations of how the work may be performed safely
Section Two: Further Reading
Types of Delivery Systems
The delivery of electric energy from the generating plant to the consumer may consist of
several more or less distinct parts that are nevertheless somewhat interrelated. The part
considered 'distribution', i.e., from the bulk supply substation to the meter at the
consumer's premises, can be conveniently divided into two subdivisions:
1. Primary distribution, which carries the load at higher than utilization
voltages from the substation (or other source) to the point where the
voltage is stepped down to the value at which the energy is utilized by the
2. Secondary distribution, which includes that part of the system operating at
utilization voltages, up to the meter at the consumer's premises.
Primary Distribution
Primary distribution systems include three basic types:
Radial Systems. The radial-type system is the simplest and the one most
commonly used. It comprises separate feeders or circuits 'radiating' out of the
station or source, each feeder usually serving a given area. The feeder may
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be considered as consisting of a main or t run k portion from which there radiate
spurs or laterals to which distribution transformers are connected, as illustrated in
Figure 4-3.
The spurs or laterals are usually connected to the primary main through fuses, so
that a fault on the lateral will not cause an interruption to the entire feeder. Should
the fuse fail to clear the line, or should a fault develop on the feeder main, the circuit
breaker back at the substation or source will open and the entire feeder will be
deenergized .
To hold down the extent and duration of interruptions, provisions are made to
sectionalize the feeder so that unfaulted portions may be reenergized as quickly as
practical. To maximize such reenergization, emergency ties to adjacent feeders are
incorporated in the design and construction; thus each part of a feeder not in trouble
can be tied to an adjacent feeder. Often spare capacity is provided for in the feeders
to prevent overload when parts of an adjacent feeder in trouble are connected to them.
Figure 4-4. Primary Feeder Schematic Diagram Showing Trunk or Main Feeds
and Laterals
or Squts .
enough diversity between loads on adjacent feeders to require no extra capacity
to be installed for these emergencies.
Loop Systems. Another means of restricting the duration of interruption
employs feeders designed as loops, which essentially provide a two-way primary
feed for critical consumers. Here, should the supply from one direction fail, the
entire load of the feeder may be carried from the other end, but sufficient spare
capacity must be provided in the feeder. This type of system may be operated
with the loop normally open or with the loop normally closed.
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Primary Network Systems. Although economic studies indicated that under
some conditions the primary network may be less expensive and more reliable than
some variations of the radial system, relatively few primary network systems
have been pu t into actual operation and only a few still remain in service.
This system is formed by lying together primary mains ordinarily found in radial
systems to form a mesh or grid. The grid is supplied by a number of power
transformers supplied in tu r n from subtransmission and transmission lines at higher
voltages. A circuit breaker between the transformer and grid, controlled by
reverse-current and automatic reclosing relays, protects the primary network from
feeding fault current through the transformer when faults occur on the supply
subtransmission lines. Faults on sections of the primaries constituting the grid are
isolated by circuit breakers and fuses. See Figure4-4.
This type of system eliminates the conventional substation and long primary
trunk feeders, replacing them with a greater number of ‘unit’ substations strategically
placed throughout the network. The additional sites
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necessary are often difficult to obtain. Moreover, difficulty is experienced in
maintaining proper operation of the voltage regulators (where they exist) on the
primary feeders when interconnected.
Secondary Distribution
Secondary distribution systems operate at relatively low utilization voltages and, like
primary systems, involve considerations of service reliability and voltage regulation.
The secondary system may be of four general types:
Individual Transformers-Single Service. Individual-transformer service
is applicable to certain loads that are more or less isolated, such as in rural
areas where consumers are far apart and long secondary mains are
impractical, or where a particular consumer has an
extraordinary large or unusual load even though
situated among a number of ordinary consumers.
In this type of system, the cost of the
several transformers and the sum of power losses
in the units may be greater (for comparative
purposes) than those for one transformer supplying a group of consumers from its associated
secondary main. The diversity among consumers'
loads and demands permits a transformer of smaller capacity than the capacity of the
sum of the individual transformers to be installed. On the other hand, the cost and losses
in the secondary main are obviated, as is also the voltage drop in the main. Where low
voltage may be undesirable for a particular consumer, it may be well to apply this type
of service to the one consumer. Refer to Figure 4-6.
Figure 4-6. Common-Secondary-Main Supply.
Common Secondary Main. Perhaps the most common type of secondary system
in use employs a common secondary main. It takes advantage of
diversity between consumers' loads and demands, as indicated above. Moreover, the
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larger transformer can accommodate starting currents of motors with less resulting
voltage dip than would be the case with small individual transformers. Sec Figure
Banked Secondaries. The secondary system employing banked second- aries is
not very commonly used, although such installations exist and are usually limited to
overhead systems.
This type of system may be viewed as a single-feeder low-voltage network, and the
secondary may be a long section or grid to which the transformers arc connected.
Fuses or automatic circuit breakers located between the transformer and secondary
main serve to clear the transformer from the bank in case of failure of the transformer.
Fuses may also be placed in the secondary main between transformer banks. See
Figure 4-7.
Some advantages claimed for this type of system include uninterrupted service,
though perhaps with a reduction in voltage, should a transformer fail; better
distribution of load among transformers; better normal voltage conditions resulting
from such load distribution; an ability to accommodate load increases by changing
only one or some of the transformers, or by installing a new transformer at some
intermediate location without disturbing the existing arrangement; the possibility
that diversity between demands on adjacent transformers will reduce the total
transformer load; more capacity available for inrush currents that may cause flicker;
and more capacity as well to burn secondary faults clear.
Some disadvantages associated with this type of system are as follows: should
one transformer fail, the additional loads imposed on adjacent units may cause them
to fail, and in turn their loads would cause still other transformers to fail (this is
known as cascading).
Secondary Networks. Secondary networks at present provide the highest degree
of service reliability and serve areas of high load density, where revenues justify their
cost and where this kind of reliability is imperative. In some instances, a single
consumer may be supplied from this type of system by what are known as spot
In general, the secondary network is created by connecting together the secondary
mains fed from transformers supplied by two or more primary feeders. Automatically
operated circuit breakers in the secondary connection between the transformer and the
secondary mains, known as network protectors, serve to disconnect the transformer
from the network when its primary feeder is deenergized; this prevents a back feed
from the secondary
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into the primary feeder. This is especially important for safety when the primary feeder
is deenergized from fault or other cause. The circuit breaker or protector is backed up by
a fuse so that, should the protector fail to operate, the fuse will blow and disconnect
Figure 4-7. Banked Secondary Supply.
The number of primary feeders supplying a network is very important.
With only two feeders, only one feeder may be out of service at a time, and
there must be sufficient spare transformer capacity available so as not to
overload the units remaining in service; therefore this type of network is
sometimes referred to as a single-contingency network.
Most networks are supplied from three or more primary feeders, where the
network can operate with the loss of two feeders and the spare transformer capacity
can be proportionately less. These are referred to as second-contingency networks.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. A radial distribution system .......... .
a. consists of a source from which mains and laterals radiate
b. consists of the trunks from which laterals equipped with transformers
c. includes that part of the system that employs feeders designed as
open loops.
d. includes that part of the system that employs feeders designed as
closed loops.
2. The circuit breaker at the source deenergizes the entire feeder ....…... .
a. if a fault develops on the feeder main
b. as soon as the fuses installed at the lateral and main intersections
begin to work
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if any of the fuses installed at each lateral and main intersection fails
to operate
d. both a and c
3. In order to minimize the duration of interruptions, .......... .
a. feeders are sectionalized
b. adjacent feeders are provided with emergency ties
c. spare capacity is provided in the feeder
d. all of the above
4. It is true that .......... .
a. primary network systems are more popular than radial systems
b. voltage regulators installed on the primary network feeders do not
always operate properly
c. primary network systems are usually less reliable than radial systems
d. circuit breakers between the grid and the transformers do not
function very well
5. Individual-transformer service is applied .......... .
a. where an individual consumer has an extraordinary large load
b. where an individual consumer is not distant from the source
c. to increase power reliability in thinly-populated areas
d. to reduce the sum of power losses in the units
6. One crucial disadvantage of a low-voltage network is that .......... .
a. distribution of load among transformers is not uniform
b. additional transformers are difficult to be installed
c. if one transformer fails, other transformers may also fail
d. if flickers appear, they cannot be prevented
7. A network protector in a secondary network is employed .......... .
a. to serve as a path between the transformer and the secondary mains
b. to deenergize the primary feeder and disconnect the transformer
from the network
c. to blow the fuse and disconnect the transformer from the secondary
c. to prevent a back feed from the secondary mains into the primary
B. Write the answers to the following questions.
1. What are the two subdivisions of the distribution system?
2. What does the primary system include?
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3. What is the function of the fuses that connect the laterals to the mains?
4. How does the loop system restrict the duration of interruption?
5. What part does the spare capacity provided in the loop system play?
6. What is the function of a circuit breaker installed between the transformer and the grid in a primary network system?
7. How are faults on the sections of the primaries isolated?
8. What replaces conventional substations in the primary network system?
9. What type of secondary distribution system is applied where long secondary
mains are impractical?
10. What is the advantage of a large transformer over a small one?
11. What is cascading?
12. What characteristics should the transformers employed in a low-voltage network
13. How does a low-voltage network handle the diversity of demands?
14. What are the single-contingency and second-contingency networks?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Each item selected, and its place in the distribution system, must be viewed in light of
its possible failure or malfunction for whatever reason. The design of the distribution
system, therefore, should take into consideration the method of maintaining each of
the several elements making up the distribution system.
The Distribution System. In general, the safest means is to be able to deenergize
and ground the particular item requiring maintenance, preferably without affecting the
remainder of the circuit. Circuits, both primary and secondary, are arranged so that small
sections may be deenergized by-interconnecting the remaining portions to other
sources, by means of some sort of switches. Smaller pieces may be deenergized by
means of hot-line or live-line clamps.
Where decncrgization is not practical, work may be carried o ut by insulating the
worker. This is accomplished by protective gear, such as rubber gloves, sleeves,
blankets, line
hose, insulator
hoods, and other similar devices
Another method insulates workers from ground by having them work from insulated
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platforms or insulated buckets mounted on line trucks. Still another means involves the
handling of the energized facilities by tools having sufficient insulation properties
which, properly handled, enable the worker to accomplish required maintenance by
essentially a remote operation of the tools; this is referred to as live-line, or hot-line,
maintenance. To facilitate such operations, appropriate details and modifications are
included in the design of distribution systems, e.g., wider spacing of conductors;
hot-line ties that hold conductors to the insulators; 'unnecessary' extensions of primary
and secondary mains so that mains butt each other, permitting their temporary
connection during contingencies by means of jumpers and the arrangement of
terminals at substations to accommodate portable substations.
B. Find the Persian equivalents of the following terms and expressions
and write them in the spaces provided.
1. circuit fault
2. compensate
3. emanate
4. facility
5. flicker
6. inrush current
7. 1ateral
8. limb
9. malfunction
10. obviate
11. overload
12. premise
13. radial system
14. regulator
15. sectionalizing fuse
16. service drop
17. short circuit
18. single-contingency network
19. spare capacity
20. spot network
21. spur
22. substation bus
23. sustain
24. trunk
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Section One: Reading Comprehension
Protective Devices
For the distribution system to function satisfactorily, faults on any part of it must be
isolated or disconnected from the rest of the system as quickly as possible; indeed, if
possible, they should be prevented from happening. The principal devices to accomplish
this include fuses, automatic sectionalizers, reclosers, circuit breakers, and lightning or
surge arresters. Success, however, depends on their coordination so that their operations
do not conflict with each other. Figure 5-1 indicates where these devices are connected
on the system.
Time-Current Characteristic. A fuse consists basically of a metallic element that melts
when ‘excessive’ current flows through it. The magnitude of the excessive current will
vary inversely with its duration. This time-current characteristic is determined not
only by the type of metal used and its dimensions (including its configuration), but also
on the type of its enclosure and holder. The latter not only affect the melting time, but in
addition, affect the arc clearing time. The clearing time of the fuse, then, is the sum of
the melting time and the arc clearing time. Refer to Figures 5-1 and 5-2. Note that for
curve b in Figure 5-3, the clearing time for a certain value of current
is less than for curve a; the fuse with the characteristic b is therefore referred
to as a ‘fast’ fuse, compared with the fuse of curve a.
Fuses are rated in terms of voltage, normal current-carrying ability, and
interruption characteristics usually shown by time-current curves. Each curve
actually represents a band between a minimum and a maximum clearing time
for a particular fuse.
Fuse Coordination. The number, rating, and type of the interrupting devices
shown in Figure 5-1 depend on the system voltage, normal current, maximum fault
current, the sections and equipment connected to them, and other local conditions. The
devices are usually located at branch intersections and at other key points, When two or
more such devices are employed in a circuit, they will be coordinated so that only the
faulted portion will be deenergized.
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Figure 5-1. Radial Primary Feeder Showing Location of Protective Devices.
Repeater Fuses
Line fuses are sometimes installed in groups
of two or three ( per phase ), known as
repeater fuses, having a time delay between
each two fuse units. When a fault occurs, the
first fuse will blow and the second fuse
will be mechanically placed in the circuit by
the opening of the first; if the fault Persists, the
second fuse will blow; if
There a third fuse,
the process is
repeated. If the fault is permanent, all
of the fuses will blow and the faulted part of the circuit will be deenergized. new fuses
must be installed to restore the line to normal.
Where capacitors are applied to feeders for power factor correction,
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fuses chosen to protect the line from the bank (and vice versa) must also coordinate
with sectionalizing and other devices in the circuit back to the source .
Transformer Fuses
Fuses on the primary side of distribution transformers serve to disconnect the
transformer from the circuit not only in the event of a fault in the transformer or on
the secondary, but also when the normal load on the transformer becomes so high
that failure is imminent. Fuses on the secondary side protect the transformer from
faults or overloads on the secondary circuit it serves.
The characteristics of a primary fuse are a compromise between protection
from a fault and protection from overload, yet the fuse also has to coordinate with
other fuses on the line. One attempt at a solution is the completely self-protected
(CSP) transformer, in which the primary fuse, with characteristics based only on
protection against fault, is situated within the transformer tank (and, to
differentiate, is called a link) while overload protection is accomplished by
low-voltage circuit breakers (instead of fuses) on the secondary side of the
transformer that are also situated within the tank. The circuit breakers, once open,
however, must be reclosed manually.
of the protectors on low-voltage
Fuses are provided on the line side
secondary networks. These are backup
protection in the event the protector
fails to open during back feed from the
network into the primary when it is faulted or
deliberately grounded.
Secondary fuses, known as limi-ters,
are also provided at the juncture of secondary
mains to isolate faulted sections of the
secondary mains and to prevent the spread of
burning in conductors (usually in cables)
where sufficient fault current does not exist
to burn them clear in a small portion of the
Part I. Comprehension Exercises
A. Put T for true and "F" for false statements. Justify your answers.
....... 1. Protective devices connected on a distribution system cause the
system to function satisfactorily.
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…… 2. The greater the excessive current flowing through a fuse, the
longer the melting time.
…… 3. Different fuses with different rates of voltage, current-carrying
ability, and interruption characteristics can be produced.
…… 4. The interrupting devices may be located anywhere in a distribution
…… 5. Limners, employed at the juncture of the secondary mains, isolate
their faulted sections.
B. Choose a, b, c, or d which best completes each item.
1. The clearing time of a fast fuse is .......... .
a. comparatively low
b. comparatively high
c. equal to its arcing lime d. equal to its melting time
2. According to the passage, .......... .
a. the minimum and the maximum of the clearing time for a certain
value of current cannot be evaluated
b. the minimum and the maximum of the clearing time of a fuse are
always constant
c. protective devices in a circuit must be coordinated so that their
operations do not conflict with each other
d. protective devices in a circuit must be adjusted to deficiencies
resulting from manufacturing problems
3. Repeater fuses .......... .
a. are installed in groups with a time delay between each two fuse units
b. are installed in series to restore the equipment to normal
c. may act as capacitors for power factor correction
d. may all operate simultaneously to prevent deenergization
4. Fuses on the primary side of distribution transformers .......... .
a. will not protect the transformer if a fault occurs on the secondary
b. will not protect the transformer if a fault occurs in the transformer
c. will protect the fuses on the secondary side if a permanent fault
d. will disconnect the transformer from the circuit if it is seriously
5. It is true that ……. .
a. a transformer fuse is basically designed for fault protection
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b. a transformer circuit breaker is manually opened and closed
c. a self-protected transformer is equipped with links and circuit
d. a self-protected transformer is fully automatic
C. Answer the following questions orally.
1. What kinds of protective devices are used in a distribution system?
2. What does the time-current characteristic depend on?
3. How important are the enclosure and the holder of the fuse?
4. What do the number, rating, and type of the interrupting devices
applied in a circuit depend on?
5. What can a primary fuse do?
6. How are fuses used as backup protection on low-voltage secondary
Part IL Language Practice
A. Choose a, b, c, or d which best completes each item.
1. Devices called .......... are designed to open when a fault occurs on that
part of the main in which they are connected.
a. regulators
b. fuses
c. reclosers
d. relays
2. Some .......... are designed to open in air, with special provisions for
handling the arc that follows when the contacts are opened.
a. fuses
b. arresters
c. line sectionalizers
d. circuit breakers
3. In liquid-filled construction, .......... link is enclosed in a tube that is
filled with a fire-extinguishing fluid such as carbon tetrachloride.
a. the fuse
b. the switch
c. the fault-counting relay
d. the circuit breaker
4. Switches are installed in the main of the feeder, enabling the main to be
sectionalized, isolating the fault between two switches or other ……..
a. connecting
b. sectionalizing
c. energizing
d. feeding
5. Like the secondary circuit, the design of the primary .......... is based on
the maximum voltage variation permissible at the farthest consumer.
a. recloser
b. limiter
c. conductor
d. feeder
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B. Fill in the blanks with the appropriate form of the words
1. Energize
a. In a dependent time delay relay, the time delay varies with the value
of the .......... quantity.
b. A relay is said to ‘pick up’ when it changes from the unenergized
position to the .......... position.
c. Relay functioning time is the time between .......... and operation or
between deenergization and release.
d. The elapsed time after the coil has been .......... to the time required
to seat the armature of the relay is the relay seating time.
2. Blow
a. In a transformer circuit, the fuse is chosen so as to carry the inrush
transient without ......... .
b. If a fault persists in a circuit, the fuses may .......... and the faulted
part of the circuit will be deenergized.
c. A.......... blade is an active element of a fan.
d. In liquid-filled fuses, a spring is used to hold the fuse under tension
so that, when it ........... the resultant arc is quickly lengthened and
quenched in the fluid; the gas formed is inert and helps in .......... out
the arc.
3. Melt
a. Melting-speed ratio refers to the ratio of the current magnitudes
required to .......... the current-responsive element at two specified
melting times.
b. The time required for overcurrent to severe the current-responsive
element is known as the .......... time.
c. In a sand-filled fuse, the heat and gases generated when the fuse
........... are
absorbed by the sand, which tends to squelch the arc.
4. Interrupt
a. An interrupted continuous wave is a wave that is .......... at a constant
audio-frequency rate.
b. An .......... is designed to interrupt specified currents under specified
c. The loads that can be .......... in the event of a capacity deficiency on
the supplying system are .......... loads.
d. There will be the loss of service to one or more consumers or other
facilities if.......... occurs.
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5. Inverse
a. The dependent lime delay relay, known as inverse time delay relay,
has an operating lime which is an .......... function of the electrical
characteristic quantity.
b. The inverse lime delay relay with definite minimum (I.D.M.T) is a
relay in which the time delay varies .......... with the characteristic
quantity up to a certain value, after which the time delay becomes
substantially independent.
C. Fill in the blanks with the following words.
The number and sources of supply subtransmission .......... to the
distribution substation will depend not only on the .......... to be served, but
also on the degree of .......... reliability sought. Some rural substations may be
.......... from only one subtransmission feeder,
.......... serving urban and
suburban areas have a/an .......... of two supply feeders and may have ..........
more. Each additional incoming feeder, however, .......... to the bus and
switching requirements, including .......... devices for their protection, all of
.......... add to costs.
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. The higher the distribution voltage, the farther apart substations may
be located, but they also become larger in capacity and in the number
of customers served.
b. In general, it should be situated as close to the load center to be served
as practical.
c. The difficulty in obtaining substation sites is an important factor in
selecting the distribution voltage, both in original designs and in later
d. Thus, the problem of the number and location of distribution substations involves not only the study of transmission and subtransmission
designs, but more emphasis on service reliability and consideration of
additional costs that may be justified.
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e. This implies that all loads can be served without undue voltage
regulation, including future loads that can be expected in a reasonable
of time.
f. Perhaps the first consideration regarding a distribution substation is its
Section Two: Further Reading
Automatic Line Sectionalizers
Automatic line sectionalizers are connected on the distribution feeder in series with
line and sectionalizing fuses; they are also in series with and electrically farther from
the source than reclosers or circuit breakers with reclosing cycles. These devices are
decreasing in usage, but many exist on distribution systems.
When a fault occurs on the circuit beyond the sectionalizer, the fault current
initiates a fault-counting relay that is coordinated with the characteristics of the fuses
and other devices. Each time the circuit is deenergized (from reclosers or circuit
breakers), the relay moves toward the trip position; just before the final operation that
will lock out the recloser or circuit breaker if the fault persists, the sectionalizer will
trip (while no fault current is flawing) and open the circuit at that point, removing the
fault and permitting the circuit breaker or recloser to close and reset into its normal
position; service is thus restored to the rest of the circuit up to the location of the
sectionalizer. If the fault is of a temporary nature and is cleared before the reclosing
devices complete their operations, the sectionalizer will reset to its position after the
circuit is reenergized.
Sectionalizers are rated on continuous current-carrying capacity, min-tripping and
counting current, and maximum momentary fault current,
as well as for maximum
system voltage, load-break current, and impulse voltage or basic insulation level
More than one sectionalizer can be connected in series w i t h a reclosing
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device. The sectionalizer nearest the reclosing device can be set to operate
after (say)
three operations while the more remote one is set for (say) two
such operations.
Sectionalizers are relatively low-cost devices; they are not required to interrupt fault
current although fault current flows through them. They may
be operated manually and
are considered the same as load-break switches.
Reclosers are essentially circuit breakers of lower capacity, both as to normal current and
interrupting duty. They are usually installed on major branches of distribution feeders in
series with other sectionalizing devices; they perform the same function as repeater fuses
connected in the circuit or circuit breakers
at the substation.
Reclosers are designed to remain open, or ‘locked out’, after a selected sequence of
tripping operations. A fault will trip the recloser; if the fault is temporary in nature and no
longer exists, the next tripping operation does not take place and the recloser returns to its
normally closed position, ready for another incident. If the fault persists, the recloser will
close and the operation will be repeated until the recloser locks out. The reclosers are
usually set for three automatic reclosing operations before locking out; the first operation is
usually ‘instantaneous’, i.e., occurring as quickly as the breaker contacts can open with no
time delay; the second and third operations have time delays inserted, that for the second
tripping smaller than that for the third; a fourth tripping will result in the recloser’s
remaining open until it is automatically or manually restored to normal, ready for the next
Reclosers can operate on one or more time-current characteristic curves.
reclosing characteristics of the recloser for each operation are coor-dinated with those
of the fuses at the coordinating points in the circuit and
with those of the relays
controlling the circuit breaker at the substation.
Circuit Breakers-Relays
Where the fault current is beyond the ability of a fuse or recloser to interrupt it safely,
or where repealed operation within a short period of lime makes it more economical, a
circuit breaker is used. The circuit breakers must not only interrupt the normal load
current, but must be mechanically able to withstand the forces resulting from the large
magnetic fields created by the fault current flowing through them. Since the field will
depend on the magnitude of the fault current, which in turn also depends on the
voltage of the circuit, the
stresses that must be accommodated depend on both of these values. Their time-current
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characteristics, however, are dependent on the protective relays associated with them
and must be coordinated with those of down-line reclosers, fuses, and other protective
Overcurrent Relays
Overcurrent relays close their contacts to actuate the circuit that causes the circuit
breaker to open or close when the current flowing in them reaches a predetermined value.
Instantaneous. Without time delay deliberately added, the relay will close its
contacts ‘instantaneously’, i.e., in a relatively short time, in the nature of 0.5 to perhaps
20 cycles. To prevent frequent operation of the breaker from transient, nonpersistent
conditions, undesirably high settings may be applied to the relay.
Inverse Time. The operation of the relay may be made to vary approximately
inversely with the magnitude of the current. The current setting may be varied and time
delay introduced by varying the restraint on the movable element of the relay. Greater
selectivity between relays and fuses in the circuit may thus be obtained.
Definite Time. A definite time delay can be introduced before the relay begins to
operate, allowing greater selectivity to be achieved. This feature is often added to the
inverse-time characteristic beyond a certain value of current after which the relay
operation is completed after the fixed time delay. This inverse definite minimum time
feature is employed in most Overcurrent relay applications.
The distribution circuit may be sectionalized with reclosers, automatic
sectionalizers, and fuses, at which points faults may be isolated without affecting the
entire circuit; fuses are also provided on the primary side of distribution transformers.
The definite time characteristic of the relay associated with the circuit breakers at the
substation is coordinated with the characteristics of reclosers and fuses on the distribution
Directional Relays
Directional relays are essentially Overcurrent relays to which an element similar to a
wattmeter is added, both sets of contacts being in series. The Overcurrent element will
operate to close its contacts regardless of the direction of flow of power in the line; the
wattmeter element will tend to turn in one direction under normal flow of power and
in the reverse direction
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when power flows in the opposite direction. Hence, both sets of contacts must be closed
and power flowing in a given direction before the relay will operate. Both elements may be
combined into one so that only a single set of contacts is required.
This type of relay is used in primary or secondary network operations to open the
protectors to prevent current from the network from energizing the high side of the
transformers and their supply feeder during contingencies.
Differential Relays
Differential relays operate on the difference between the current entering the line or
equipment being protected and the current leaving it. As long as the incoming current
and the outgoing current are essentially equal, the relay will not operate. A fault within
the line or equipment, however, will disturb this equilibrium, and the relay will operate
to trip the supply circuit breaker or breakers on both sides of the line or equipment being
protected. This type of relay is used to protect buses, transformers, and regulators at the
substation. Since the voltages at which these operate may be high, current transformers
installed on both sides of the equipment, with proper ratios in the case of transformers,
supply the currents to the relay. Refer to Figure 5-4.
Figure 5-4. One-Line Diagram of Current-Differential Protection.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. An automatic line sectionalizer is .......... .
a. a self-controlled circuit opening device that locks out other protective devices and allows the circuit to be reenergized
b. a self-controlled circuit opening device that removes the fault and
allows the circuit to be reenergized
c. employed in parallel with other sectionalizers to prevent deenergization of the circuit
d. employed in a circuit to open the reclosers and prevent deener- gization
of the circuit
2. Reclosers are designed so that they will .......... after a selected sequence
of tripping operations.
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a. not react 10 the fault
b. not act as fuses
c. remain open or locked out
d. allow the fault to pass through them
3. It is true that reclosers .......... .
a. are set for three simultaneous operations
b. are set for indefinite automatic operations
c. can operate on only one time-current characteristic
d. can operate on one or more time-current characteristic curves
4. As we understand from the text, .......... .
a. reclosers are not designed to withstand the forces resulting from
severe faults
b. reclosers are designed to do the same function as circuit breakers in
the circuit
c. circuit breakers are not able to withstand the stresses resulting from
large magnetic fields
d. circuit breakers are independent from relays associated with them
and other protective devices in the circuit
5. A definite-time relay indicates that .......... .
a. a delay is purposely introduced in action regardless of the magnitude
of the quantity that causes the action
b. a delay introduced in action causes the relay to operate inversely
with the magnitude of the current
c. the relay may close its contacts instantaneously regardless of the
addition of a time delay
d. the relay may be based on a time delay introduced by varying the
restraint on its movable element
6. A differential relay is applied to the circuit .......... .
a. to respond to the difference between the current entering the
protected equipment and the current leaving it
b. to respond to the difference between the high voltage currents
flowing through the protected apparatus
c. to stabilize the state of equilibrium between the current entering the
protected line and the current leaving it
d. to stabilize the circuit breakers on both sides of the protected
equipment with enough energy
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B. Write the answers to the following questions.
1. How are automatic line sectionalizers employed in a circuit?
2. How are sectionalizers rated?
3. How does a sectionalizer close to a reclosing device function compared
with one which is far from the reclosing device?
4. Where are reclosers used as repeater fuses?
5. How are reclosers compared with circuit breakers?
6. What is the function of an overcurrent relay?
7. What are the different types of overcurrent relays?
8. What part does a directional relay play in a distribution system?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Surge or Lightning Arresters
The function of a surge or lightning arrester is to limit the voltage stresses on the
insulation of the equipment being protected by permitting surges in voltage to drain to
ground before damage occurs. The surges in voltage generally are caused by lightning
(either by direct stroke or by induction from a nearby stroke) or by switching.
Arresters consist of two basic components: a spark gap and a nonlinear resistance
element (for a valve type) or an expulsion chamber (for an expulsion type). When a
surge occurs, the spark gap breaks down or sparks over, and permits current to flow
through the resistance (or chamber) element to ground. Since the arrester at this point
presents a low-impedance path, a large current, referred to as 60-cyde follow current,
flows through the arrester. The nonlinear resistance, at the higher voltages, will tend to
restrict this current and eventually cause it to cease to flow; here, the magnitude of the
follow current is independent of the system capacity. The expulsion chamber will
confine the arc, build up pressures that eventually blow out the arc, and cause the
follow current to cease to flow; here, the follow current is a function of the system
capacity and the expulsion chamber must be suitably designed. After each such
operation, the arrester must be capable of repeating this operating cycle.
B. Find the Persian equivalents of the following terms and
write them in the spaces provided.
expressions and
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2. contingency
3. deficiency
4. elapse
5. expulsion chamber
6. fault-counting relay
7. imminent
8. instantaneous
9. lighting arrester
l0. link
11. permanent
13. quench
15. repeater fuse
16. restraint
17. squelch
18. tension
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Section One: Reading Comprehension
Relationship Between Voltage and
This is a long and interesting story. It is the heart of electronics. Crudely speaking, the
name of the game is to make and use gadgets that have interesting and useful I versus V
characteristics. Resistors (I simply propor tional to V ), capacitors (I proportional to rate
of change of V ), diodes (I
only flows in one direction), thermistors
(temperature-dependent resistor), photoresistors (light-dependent resistor), strain gauges
(strain-dependent resistor), etc., are examples. We will gradually get into some of these
devices; for now, we will start with
the most mundane and most widely
used circuit element, the resistor
(Figure 6-1).
It is an interesting fact that
the current through a metallic conductor (or other partially conducting material) is
proportional to the voltage across it. (In the case of wire conductors used in circuits, we
usually choose a thick enough gauge of wire so that these Voltage drops' will be
negligible.) This is by no means a universal law for all objects. For instance, the current
through a neon bulb is a highly nonlinear function of the applied voltage (it is zero up to a
critical voltage, at which point it rises dramatically). The same goes for a variety of
interesting special devices-diodes, transistors, light bulbs, etc.
A resistor is made out of some conducting stuff (carbon, or a thin metal or carbon
film, or wire of poor conductivity), with a wire coming out each end. It is characterized
by its resistance:
R is in ohms for V in volts and I in amps. This is known as Ohm's law. Typical resistors of
the most frequently used type (carbon composition) come in values from 1 ohm (1  )
to about 22 megohms (22 M  ). Resistors are also characterized by how much power
they can safely dissipate (the most
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commonly used ones are rated at 1 4 or 1 2 watt and by other parameters such
as tolerance (accuracy), temperature coefficient, noise, voltage coefficient (the
extent to which R depends on applied V ) , stability with time, inductance, etc.
Capacitors (Figure 6-2) are devices that might be considered simply
frequency-dependent resistors. They allow you to make frequency-dependent
voltage dividers, for instance. For some applications (bypass, coupling) this is
almost all you need to know, but for other applications (filtering, energy
storage, resonant circuits) a deeper understanding is needed. For example,
capacitors cannot dissipate power,
even though current can flow
through them, because the voltage
and current are 90° out of phase.
A capacitor ( the old-fashioned
name was condenser) is a device that has two wires sticking out of it and has
the property
Q = CV
A capacitor of C farads with V volts across its terminals will contain Q coulombs of
stored charge.
Taking the derivative, you get
I= C
Capacitors come in an amazing variety of shapes and sizes. The basic construction is
simply two conductors near each other; in fact, the simplest capacitors are just that. For
greater capacitance, you need more area and closer spacing; the usual approach is to
plate some conductor onto a thin insulating material (called a dielectric), for instance,
aluminized Mylar film rolled up into a small cylindrical configuration. Other popular
types are thin ceramic wafers, metal foils with oxide insulators (electrolytics), and metallized
mica. Each of these types has unique properties. In general, ceramic and Mylar types
are used for most noncritical circuit applications; tantalum capacitors are used where
greater capacitance is needed, and electrolytics are used for power supply filtering.
Inductors (Figure 6-3) are
closely related to capacitors; the
rate of current change in an
inductor depends on the voltage
applie across it, whereas the rate
of voltage change in a capacitor depends on the current through it. The
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defining equation for an inductor is
where L is called the inductance and is measured in henrys (or mH ,  H, etc.). Putting
a voltage across an inductor causes the current to rise as a ramp (for a capacitor,
supplying a constant current causes the voltage to rise as a ramp).
The symbol for an inductor looks like a coil of wire; that is because, in its simplest
form, that is all it is. Variations include coils wound on various core materials, the most
popular being iron (or iron alloys, laminations, or powder) and ferrite, a black,
nonconductive, brittle magnetic material. These are all ploys to multiply the inductance
of a given coil by the ‘permeability’ of the core material. The core may be in the shape
of a rod, a toroid (doughnut), or even more bizarre shapes, such as a ‘pot core’ (which
has to be seen to be understood; the best description we can think of is a doughnut mold
split in half, if doughnuts were made in molds).
Inductors find heavy use in radio frequency (RF) circuits, serving as RF ‘chokes’
and as parts of tuned circuits. A pair of closely coupled inductors form the interesting
object known as a transformer.
V L
Part I. Comprehension Exercises
A. Put T for true and "F" for false statements. Justify your answers.
........ 1. The relationship between voltage and current is a crucial point in
........ 2. Resistors and capacitors provide for useful current versus voltage.
........ 3. Thermistors are light sensitive devices.
........ 4. Capacitors are of different types.
........ 5. Each capacitor has various applications.
........ 6. In order to produce coils with different rates of inductance, various
core materials are employed.
B. Choose a, b, c, or d which best completes each item.
1. Ohm’s law states that the current in a circuit is .......... .
a. inversely proportional to the resistance of the circuit and is directly
proportional to the electromotive force in the circuit
b. directly proportional to the resistance of the circuit and is inversely
proportional to the electromotive force in the circuit
c. directly proportional to the resistance and the electromotive force in
the circuit
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d. inversely proportional to the resistance and the electromotive force
in the circuit
2. It is true that resistors .......... .
a. separate signals
b. generate waves
c. dissipate power
d. store energy
3. The current through a capacitor is .......... .
a. independent of frequencies
b. proportional to the voltage
c. independent of the variations of voltage
d. proportional to the rate of change of voltage
4. It is true that .......... .
a. resistors may be used in bypass applications
b. capacitors may be used for filtering
c. resistors are identical with condensers
d. capacitors are identical with thermistors
5. If you change the voltage across a farad by one volt per second, you are
……. .
a. supplying an ampere
c. increasing the voltage
b. supplying a farad
d. increasing the current
6. We may deduce from the text that 1 volt across 1 henry produces a
curent that .......... .
a. decreases at 1 amp per second
b. increases at 1 amp per second
c. is constant up to a critical point
d. is zero up to a critical point
C. Answer the following questions orally.
1. What is the function of a diode?
2. What is a resistor made up of?
3. How are capacitors basically made?
4. Which part of a capacitor is called a dielectric?
5. Why does aluminized Mylar film act as a capacitor?
6. How do you describe a pot core used in an inductor?
7. What are some applications of inductors?
Part II. Language Practic
A. Choose a, b, c, or d which best completes each item.
1- A voltage or current that varies at a constant rate is referred to as
a. ramp
b. a rise
c. a drop
d. a tap
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2. A .......... is an electron device that makes use of the change of resistivity
of a semiconductor with change in temperature.
a. thermocouple
b. thermoelement
c. thermostat
d. thermistor
3. A .......... consists of two electrodes separated by a dielectric for
introducing capacitance into an electric circuit.
a. capacitor
b. resistor
c. diode
d. strain gauge
4. A .......... introduces relatively small insertion loss to waves in one or
more frequency bands and relatively large insertion loss to waves of
other frequencies.
a. bypass
b. filter
c. coupling
d. condenser
5. The property of an electric circuit by virtue of which a varying current induces
an electromotive force in that circuit or in a neighboring circuit is called .......... .
a. conductance
c. inductance
d. resistance
B. Fill in the blanks with the appropriate form of the words
1. Recognize
a. Electrical signal such as voltages exist throughout a digital system in. either one of
two ........... values and represent a binary variable equal
to 1 or 0.
b. Each digital logic family is ........... by its basic NOR or NAND.
2. Rate
a. Rated accuracy is the limit that errors will not exceed when an instrument is
used under any combination of .......... operating
b. Rate-of-rise suppressors are devices used to control the .......... of rise
of current and/or voltage to the semiconductor devices in a semiconductor power
c. The .......... of electric apparatus in general is expressed in voltamperes,
horsepower, kilowatts, or other appropriate units.
3. Resist
a. A resistor introduces .......... into an electric circuit.
b. A .......... as used in electric circuits
of operation,
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protection, or control, commonly consists of an aggregation of units.
c. When the resistivity of substance is known, the .......... of any body
composed of that substance can be calculated.
d. The resistance of a wire is directly proportional to the .......... of the
substance forming the wire.
4. Capacitor
a. A parallel circuit consisting of an inductor in parallel with a .......... is
termed a parallel resonant or parallel tuned circuit when the resultant
current taken from the supply is at its minimum value.
b. A filter consists of an arrangement of resistors and inductive and .........
c. Capacitance current or component is a reversible component of the
measured current on charge or discharge of the winding and is due
to the geometrical ........... , that is, the capacitance as measured with
altering current of power or higher frequencies.
5. Depend
a. The distinction between the two aspects of reliability, .......... and
security, is usually made when a communication channel is involved
in the relay system and noise or extraneous signals are a potential
hazard to the correct performance of the system.
b. A dependent contact is a contacting member designed to complete
any one of two or three circuits, .......... on whether a two- or
three-position device is considered.
c. An operation solely by means of directly applied manual energy is
referred to as ......... manual operation.
C. Fill in the blanks with the following words.
A series tuned circuit is .......... of a capacitor and an inductor .......... in
series. The frequency at which the .......... and inductive reactances are equal is
the ...... frequency. The reactance value of the .......... or capacitor at this
frequency .......... by the series resistance in the .......... is the Q (quality factor).
A Parallel ........... circuit consists of an inductor in .......... with a capacitor.
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D. Put the following sentences in the right order to from a
paragraph. Write the corresponding letters in the boxes
a. Inductors usually consist of many adjacent turns of insulated wire,
wound on a single support of laminated iron for low frequency
inductors and ferrite for high frequencies.
b. Resistors may be wire-wound to dissipate considerable heat, or they
may be a thin film or a composition.
c. At very high frequencies, air-core inductors are generally used.
d. Resistors, inductors and capacitors are passive circuit components.
e. Capacitors may be fixed, with solid dielectrics, or variable with usually
an air dielectric.
f. Rheostats and potentiometers are variable resistors with two and three
terminals respectively.
5 6
Section Two: Further Readin
Passive Circuit Components-Resonant
A familiar electrical component in electronic circuits is the resistor. Resistors are
circuit elements having specified values of resistance. Some resistors are made from
a long, very fine wire wound on an insulating support; and these wire-wound resistors
are generally used when it may be necessary to dissipate considerable heat.
Another common type of resistor is the thin-film resistor. This is made by
depositing a thin film of metal on a cylindrical insulating support. High resistance
values are a consequence of the thinness of the film.
A third type of resistor that has had very wide application is the
composition resistor. In this type, the resistive element is a combination of finely
divided carbon or graphite and a non-conducting inert material or filler such as
talc, with synthetic resin used as a binder. These substances are
employed in such proportions as to give the correct resistance value in the finished
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The need often arises to vary the resistance of a resistor while it is
permanently connected in a circuit. Components which enable this to be done
are known as variable resistors, and their main element is a mechanical slider
or arm which slides over the resistance element included in the circuit. A
variable resistor with only two terminals is known as a rheostat while one with three is
known as a potentiometer.
Electronic components with appreciable inductance are called inductors. They
consist of many adjacent turns of wire wound on the same support. large inductances
for use at low frequencies are obtained by winding many hundreds of turns of wire on a
core of a ferromagnetic material such as iron. When iron cores are used they are
laminated in order to reduce eddy currents. Ferrite cores, made of high resistivity
ferromagnetic materials, are used at high frequencies because their high resistance
makes eddy currents negligible. They are not used at low frequencies, however, because
their magnetic properties are not so favourable as those of iron. At very high
frequencies, air-core inductors are employed. Despite the fact that variable inductances
can be obtained by moving one portion of the winding with respect to the other, such
components are not widely used.
The third major component we may consider is the capacitor. Capacitors may be
variable or fixed. Variable capacitors most often use an air dielectric but sometimes the
dielectric can be compressed gas or a liquid. Capacitance adjustment in variable
capacitors is usually obtained by varying the effective plate area. Variable capacitors
are also termed tuning capacitors and are generally used for varying the resonance
frequency of a tuned circuit.
Fixed capacitors employ a wide range of dielectric materials and new types are
continually finding valuable applications. Among the important types are those with
solid dielectrics such as mica, plastic films, certain ceramics, paper and electrolytic
In a circuit composed of a capacitor and an inductor connected in series
with a source of alternating current in which the frequency can be varied over range, at
low frequencies the capacitive reactance of the circuit is large and the inductive
reactance is small, while at high frequencies the inductive reactance is large and the
capacitive reactance is small. Between these two extremes there is a frequency called
the resonant frequency at which the capacitive and inductive reactances are exactly
equal. As a result they completely cancel each other out and the current flow is
determined wholly by the resistance of the circuit. Under these conditions the circuit
is termed a
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Series resonant circuit. At the resonant frequency the current has its largest
Value, assuming the source voltage to be constant regardless of frequency.
The principle of resonance finds its most extensive application in radio
frequency circuits, The value of the reactance of either the inductor or th e capacitor
at the resonant frequency of a series resonant circuit divided by the series resistance in
the circuit (which is always present) is termed the quality factor, or more usually the Q
factor or merely the Q of the circuit. In particular, for a given value of inductance, a
circuit having a higher Q will have a smaller resistance and consequently will have a
higher resonant current in comparison with the off-resonance current. Consequently
(the curve of current versus frequency will be more sharply peaked and the circuit is
said to be more sharply tuned.
A parallel circuit consisting of an inductor in parallel with a capacitor is termed a
parallel resonant or parallel tuned circuit when the resultant current taken from the
supply is at its minimum value. When a variable frequency source of constant voltage
is applied to this parallel circuit there is a resonant effect similar to that in a series
circuit, although in this case the source current is smallest at the frequency for which
the inductive and capacitive reactances are equal.
A filter is a network that will permit the passage of electrical signals of a particular
frequency or band of frequencies, while offering a much greater impedance to signals
of higher or lower frequencies. Such a network may therefore be employed either to
accept or reject signals of given frequencies. It consists of an arrangement of
resistors and inductive and capacitive elements. There are three main types of
filters; low-pass, high-pass, and band-pass. A low-pass filter is one that will permit
all frequencies below a specified one-the cut-off frequency-to be transmitted with
little or no loss, it will, however, attenuate all frequencies above the cut-off
frequency. A high-pass filter is one with a cut-off frequency above which there is little
or no loss in transmission, whereas below which there is considerable attenuation. A
band-pass filter, on the other hand, will transmit a selected band of
frequencies, attenuating all those either higher or lower than the desired band.
Comprehension Exercises
A. Choose s, b, c, or d which best completes each item.
1. The first three paragraphs mainly discuss .......... wire-wound, thin-film,
and composition resistors.
a. a kind of material identical in
b. the rate of resistance offered by
c. a variety of materials used in
d. the structure and the application of
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2. The resistance values of thin-film resistors vary inversely with .......... of
the film.
a. the thickness
b. the thinness
c. the length
d. the width
3. The fifth paragraph mainly discusses various .......... .
a. inductors and their deficiencies
b. inductors and their applications
c. core materials used in inductors
d. core materials used at low frequencies
4. To vary the resonant frequency of a tuned circuit ..........
are employed.
a. variable resistors
b. variable capacitors
c. fixed capacitors
d. composition resistors
5. The eighth paragraph mainly discusses .......... .
a. the capacitive reactance at low frequencies
b. the inductive reactance at high frequencies
c. the composition of a series resonant circuit
d. the application of a series resonant circuit
6. When a variable frequency source of constant voltage is applied to a
parallel resonant circuit there is a resonant effect and .......... .
a. the source current is smallest at the frequency for which the
inductive and capacitive reactances are equal
b. the source current is largest at the frequency for which the inductive
and capacitive reactions are equal
c. the inductive and capacitive reactances cancel each other out and
produce an impedance
d. the inductive and capacitive reactances cancel each other out and
result in high resistance
7. The last paragraph mainly discusses .......... .
a. the cut-off frequency
c. the structure of a filter
d. the function of a filter
B. Write the answers to the following questions.
1. What is the application of wire-wound resistors?
2. What is a composition resistor composed of?
3. What is a variable resistor?
4.What is the difference between a rheostat and a potentiometer?
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5. Why are ferrite cores for inductors not used at low frequencies?
6. Why are iron cores in inductors laminated?
7. What are some important types of dielectric materials that fixed capacitors use?
8. What is the advantage of the resonant frequency?
9. What is the Q of a circuit?
10. What are filter networks used for?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Resistors are truly ubiquitous. There are almost as many types as there are applications.
Resistors are used in amplifiers as loads for active devices, in bias networks, and as
feedback elements. In combination with capacitors they establish time constants and act
as filters. They are used to set operating currents and signal levels. Resistors are used in
power circuits to reduce volt-ages by dissipating power, to measure currents, and to
discharge capacitors after power is removed. They are used in precision circuits to
establish currents, to provide accurate voltage ratios, and to set precise gain values. In
logic circuits they act as bus and line terminators and as ‘pull-up’ and
resistors. In high-voltage circuits they are used to measure
voltages and to equalize
leakage currents among diodes or capacitors connected in series. In radio frequency
circuits they are even used as toil forms for inductors.
Resistors are available with resistances from 0.01 ohm through 1012 ohms standard
power rating from 18 watt through 250 watts, and accuracies from 0.005% through 20%.
Resistors can be made from carbon-composition moldings, from metal films, from wire
wound on a form, or from semiconduc- tor elements similar to field-effect transistors
(FETs). But, by far , the most familiar resistor is the 1 4 or 1 2 watt carbon-composition
resistor. These are available in a standard set of values ranging from 1 ohm to 100
megohm, with twice as many values available for the 5% tolerance as for the 10% type
We prefer the Alien Bradley type AB ( 1 4 watt, 5%) resistor for general use because of
its clear marking, secure lead seating, and stable properties.
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B. Find the Persian equivalents of the following terms and
expression and write them in the space provided.
1. aggregate
2. attenuate
3. band-pass
4. binder
5. bypass
6. choke
7. coupling
8. cut-off frequency
9. dissipate
10. doughnut
11. electromotive force
12. field-effect transistor (FET)
13. inductance
14. laminate
15. parallel tune
16. passive
17. quality factor
18. rheostat
19. thermistor
20. tolerance
21. toroid
22. ubiquitous
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Section One: Reading Comprehension
Important Features of a Computer
The development process of the computer over the last 150 years has
resulted in all
computers containing a number of fundamental features:
Stored Program Control. The computer program, which is a sequence of
instructions executed one by one to perform the required data manipulation, must be
stored within the computer. This has important advantages over external program
Conditional Branching. One advantage of internal program storage is that the
next instruction to be executed need not be the next in sequence since any instruction
can be accessed as fast as any other (this is known as random access). The choice of
which instruction to execute next can therefore be based upon the result of the previous
operation or operations, giving the computer the ability to make decisions based upon the
processing it performs.
Loops and Subroutines. The ability of a program to execute a particular set of
instructions repetitively when required can produce enormous savings in the storage
needed for the program. Conditional branches can be made to loop back and repeat a
set of instructions a number of times, and commonly required subtasks within a program
can be called up from any other part of the program as required, without needing to
include the instructions of the subtask in the main program every time it is called.
Speed of Electronics. Even though the individual instructions available in
computer may be quite limited, because each instruction can be executed so last, relatively
powerful processing can be accomplished in what appears subjectively to be a very short
time. (Compare this with the speed potential of Babbage's mechanical computer.)
Cost. The cost of computing power, and particularly the cost of computer
memory is continually decreasing. It is now cheaper to store a computer instruction in
an electronic memory than to store it on a card or a piece of paper tape.
Instructions Can Modify Themselves. Although this was one of von Neumann's
original ideas embodied in the concept of stored program control, it has not been widely
used since. One reason is that it is very difficult to keep
track of what the computer is doing once the computer program has been modified from
that originally written by the programer. In microprocessors, in particular, this concept is
avoided, because the implementation of a microcontroller with a permanently fixed
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program precludes any possibility of subsequent program changes.
A Simple Computer
Figure 7-1 shows the structure of a simple computer. The computer can be split into a
number of separate components, though the components shown do not necessarily
represent the physical division between components in a real computer. For example,
the control unit and arithmetic and logic unit (ALU) are generally implemented as a
single chip, the microprocessor, in microcomputers. Similarly, the input and output unit
may be combined into a single chip in some microcomputers. Nevertheless, Figure 7-1
represents the conceptual structure of any computer from the smallest microcomputer to
the largest mainframe computer.
Figure 7-1. Conceptual Structure of a Computer.
The first requirement of any computer is a mechanism for manipulating data. This
is provided by the ALU, which can perform such functions as adding or subtracting two
numbers, performing logical operations, incrementing and decrementing numbers and left
and right shift operations. From this
very basic set of operations, more complex
processing functions can be
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generated by programming. Larger computers may provide additional more
powerful instruction (for example, multiply and divide) within the computer
instruction set.
Clearly, every computer must also include an input and an output unit. These
provide the mechanism by which the computer communicates with the outside world.
The outside world may consist of someone typing at a computer terminal and
watching the response on a screen, or it may be some equipment, for example a washing
machine, which is providing data inputs such as water temperature, water level and drum
rotation speed, and is being controlled according to the program inside the computer,
via computer outputs which switch on and off the water taps and heater, and alter the
motor speed.
The computer must include an internal memory, which serves two functions.
First, it provides storage for the computer program; second, it provides temporary
storage for data which may be generated at some point during program execution by the
ALU, but not be required until somewhat later. Such data variables must be able to be
written into the memory unit by the computer, and subsequently read back when the
data are required. The memory is organised as a one-dimensional array (or list) of words,
and each instruction or data variable occupies one or more words in the memory. Each
word is made up of a number of bits (binary digits) of storage in parallel.
The control unit of the computer controls the sequence of operations of all the
components described above, according to the instructions in the computer program.
Each instruction is fetched from the memory, and is then decoded by the control unit and
converted into a set of lower-level control signals which cause the function specified by
that instruction to be executed. When one instruction execution has been completed the
next instruction is fetched and the process of decoding and executing the instruction is
repeated. This process is repeated for every instruction in the program and only differs if
a branch instruction is encountered. In this case, the next instruction to be fetched from
the memory is taken from the part of the memory specified by the branch instruction,
rather than being the next instruction in sequence.
The final component of the computer is a clock, or fixed-frequency oscillator,
which synchronises the operation of all parts of the computer, and ensures that all
operations occur in the correct sequence. The clock frequency defines the instruction
execution speed of the computer and is constrained by the operating speed of the
semiconductor circuits which make up the computer.
part I. Comprehension Exercises
A. Put "T" for true and "F" for false statements. Justify your answers.
........ 1. At rundown access, the access time is independent of the physical
location of the data.
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........ 2. Random access increases the speed of the computer.
........ 3. At random access, the choice of which instruction to execute next
depends on the result of the following operations.
........ 4. Subsequent program changes is a new criterion.
........ 5. Microprocessors are usually programed so as to modify instructions.
........ 6. The ALU is responsible for data manipulation.
........ 7. The control unit selects, interprets, and sees to the execution of
program instructions.
B. Choose a, b, c, or d which best completes each item.
1. The computer program needed to direct data manipulation .......... .
a. does not consist of sequentially arranged instructions
b. does not have the ability to execute instructions repeatedly
c. must be stored in the internal memory
d. must be stored in the external memory
2. The computer ........... .
a. has the ability to alter the sequence of program execution
b. lacks the ability of executing a sequence of instructions repeatedly
c. cannot call up subtasks from other parts of the program
d. cannot execute a subtask any time needed to do so
3. We understand from the text that the computer .......... .
a. can produce the instructions necessary to solve problems
b. can deviate from a top-down, structured design strategy
c. does not have direct access to the internal memory
d. does not make frequent uses of the internal memory
4. Modern computers are .......... than their predecessors.
a. cheaper and less boring
b. cheaper and more powerful
c. faster but less effective
d. faster but more expensive
5. The clock is a regular time-keeping device used as a component of
……… .
a. the ALU
b. the memory
c. the arithmetic unit
d. the control unit
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6. It is true that .......... .
a. the internal memory is directly under the control of the CPU
b. the internal memory is directly under the control of the input/output
c. the ALU is responsible for choosing data from memory
d. the ALU is responsible for decoding the instructions fetched from
C. Answer the following questions orally.
1. What does a computer program consist of?
2. How do you describe random access?
3. Why is it cheaper to store instructions in an electronic memory than to
them on cards?
4. What are the basic components of a computer?
5. What is the function of the input/output units?
6. What comprises the outside world of the computer?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. By ........... the computer is able to make logical decisions based upon
the results of computation.
a. task specification
b. problem solving
c. program branching
d. program implementation
2. If .......... knows a procedure that will solve the problem, the solution
may then be coded in a selected language.
a. the programer
b. the systems analyst
c. the operator
d. the systems designer
3. The main-control program specifies the order in which each ............ in
the program will be processed.
a. main task
b. major module
c. program design
d. subroutine module
4. The primary storage section can be designed to store a fixed number of
characters or .......... in each numbered address location.
a. words
c. files
d. bases
5. Primary memory is used transiently, which means that a program, or part of it,
is kept in .......... while the program is being executed.
a. the arithmetic-logic unit b. the control unit
c. the internal storage
d. the external storage
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B. Fill in the blanks with the appropriate form of the words
1. Program
a. The first step in .......... implementation is to debug the program.
b. A .......... to be tested has generally demonstrated that it will run and
produce results.
c. Critical or lengthy operations that have been slowly carried out by
software can be converted into microprograms and fused into ..........
read-only-memory chip.
d. Applications programers code modules that have been mapped out
by the chief.......... .
2. Communicate
a. Transferring data from one location or operation to another, for use
or for further processing, is data .......... .
b. In a data .......... system, workstations and other remote I/O devices
are linked with one or more processors to capture input data and
receive output information,
c. A computer must be able to .......... with the user.
3. Develop
a. An effective approach in the programing analysis stage of program .......... is to
break down a large problem into a series of smaller tasks.
b. The second generation of computers was .......... in 1960.
c. A microcomputer has .......... into a most necessary information processing tool in
the business today.
4. Process
a. The heart of any computer system is the ........... , which consists of
primary storage, arithmetic-logic, and control elements.
b. Data .......... consists of three basic activities: capturing the input data,
manipulating the data, and managing output results.
c. Since processors of almost any size today can .......... far more data in
a second than a single set of I/O devices can supply or receive, it is
common to overlap .......... jobs.
5. Execute
a. A load module which is the result of system routines linked with an
object module is directly .......... by the computer.
b. The time necessary for .......... a program is usually indicated on the
computer print-out.
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c. Instructions and .......... cycles are synchronized by a specific number
of electric pulses produced by an electronic clock that is built in the
C. Fill in the blanks with the following words.
When a program instruction is to be ..........., the control section first
retrieves it from .......... . Next, the instruction is interpreted to determine the
.......... action. This could mean an arithmetic .......... is needed-or a compare or
a data ........... or a branch. After this is determined, the .......... part of the
control section or ALU is .......... to execute the instruction. Often this process
will .......... the use of registers.
D. Put the following sentences in the right order to form a paragraph.
Write the corresponding letters in the boxes provided.
a. Auxiliary storage devices generally provide serial access to data, in the
same way as different pieces of music are stored serially on a cassette
b. A useful general-purpose computer normally has some type of auxiliary
or back-up storage mechanism for long-term archiving of data or
c. In order to access data from such devices, a large time penalty must be
tolerated; as a result, such devices are not usually used when fast
random access is required.
d. The auxiliary storage can usually be completely detached from the
computer, and is often some type of magnetic storage medium such as a
tape or a floppy disk.
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Section Two: Further Reading
Primary Storage Components
Within any computer, both instructions and data arc stored internally in memory as
binary numbers. This is because the computer can only understand
and execute instructions coded in the binary machine language code approdate to that particular type of computer. It may appear to the user that the computer
is executing a program written in some other programing language,
such as BASIC or PASCAL. In fact, in order to execute a program written in
any other language, the program must first be converted to the computers’
own machine language. This can either be done once and for all using a
compiler, or as the program is executed using an interpreter.
Each instruction in the machine code instruction set of the computer is normally
represented by one or more words in the computer memory. Each word is represented
physically by a number of binary digits (bits) in parallel. Data are also stored as words
within the computer memory and the word-length of the computer therefore defines the
number of binary digits which the computer can manipulate simultaneously, since
arithmetic manipulations are generally performed upon one memory word at a time. The
number of bits per word is chosen by the designer of the computer, and is one measure of
the processing power of the computer. At the present time, microprocessors typically
use 8, 16 or 32 bits per word, minicomputers 16-32 bits per word and mainframe computers
32-64 or more bits per word.
Computer memory is implemented in a number of different ways. Semiconductor
memory is the dominant technology at present. Semiconductor storage elements are tiny
integrated circuits. Both the storage cell circuits and the support circuitry needed for
data writing and reading are packaged on chips of silicon. There are several
semiconductor storage technologies cur rently in use. It is not necessary to consider the
physics of these different approaches in any detail. It is enough just to mention that
faster and more
expensive bipolar semiconductor chips are often used in the
and certain other sections of the processor while slower and less
expensive chips that employ metal-oxide semiconductor (MOS) technology are usually
used in the primary storage section. These primary storage components are
referred to as random-access memory (RAM) chips because any of the locations on a
chip can be randomly selected and used to directly store and retrieve data and
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RAM chips may be classified as dynamic and static. The storage cell circuits in
dynamic RAM chips contain (1) a transistor that acts in much the same way as a
mechanical on-off light switch and (2) a capacitor that is capable of storing an electric
charge. Depending on the switching action of
the transistor, the capacitor either
contains no charge (0 bit) or does hold a charge (1 bit). Since the charge on the capacitor
tends to ‘leak off’, provision
is made to periodically ‘regenerate’ or refresh the storage
charge. A dynamic RAM chip thus provides volatile storage; that is, the data stored are
lost in the event of a power failure.
Static RAM chips are also volatile storage devices, but as long as they are
supplied with power, they need no special regenerator circuits to retain the stored data.
Since it takes more transistors and other devices to store a bit
in a static RAM, these
chips are more complicated and take up more space
for a given storage capacity than
do dynamic RAMs. Static RAMs are thus used in specialized applications, while dynamic
RAMs are used in the primary storage sections of most computers. Because of the volatile
nature of these storage elements, a backup uninterruptible power system (UPS) is
often found in larger computer installations. Personal computer users can also invest a
few hundred dollars and get a small battery-powered UPS. This device supplies current for a
period long enough for users to save data on a disk and then shut down system in an
orderly way.
Specialized Storage Elements in the Processor Unit
You know that every processor has a primary storage section that holds the active
program(s) and data being processed. In addition to this general-purpose storage section,
however, many processors also have built-in specialized
storage elements that are used
for specific processing and control purposes.
One element used during processing operations is a high-speed butter (or cache)
memory that is both faster and more expensive per character stored than primary
storage. This high-speed circuitry is used as a ‘scratch pad’ to temporarily store data and
instructions that are likely to be retrieved many times during processing. Processing speed
can thus be improved. Data may be transferred automatically between the buffer and
primary storage so that application programers are unaware of its use. Once found only
in larger systems, cache memory is now available in some of the tiny microprocessor
chips used in personal computers.
Other specialized storage elements found in many processors are used for control
purposes. The most basic computer functions are carried out by
wired circuits.
Additional circuits may then be used to combine these very
basic functions into somewhat higher-level operations (to subtract values, move data,
etc.). But it is also possible to perform these same higher-level
operations with a
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series of special programs. These programs-called microprograms because they
deal with low-level machine functions-are thus
essentially substitutes for
additional hardware.
Microprograms are typically held in the processor unit in special control storage
elements called read-only memory (ROM) chips. Unlike RAM chips,
which are
volatile, ROM chips retain stored data when the power goes off. Microprogram control
instructions that cause the machine to perform certain operations can be repeatedly read
from a ROM chip as needed, but the chip will not accept any input data or instructions
from computer users.
The most basic type of ROM chip is supplied by the computer manufac- turer as part
of the computer system, and it cannot be changed or altered by users. Such chips have
found wide application as a program storage medium
in video games and personal
computers. Of course, it is possible for a user to 'customize' a system by choosing the
machine functions that will be performed by microprograms and by then using a second
type of ROM chip. For example, critical or lengthy operations that have been slowly
carried out by software can be converted into microprograms and fused into a programable
read-only memory (PROM) chip. Once they are in a hardware form, these tasks can
usually be executed in a fraction of the time previously required.
PROM chips are supplied by computer manufacturers and custom ROM vendors. Once
operations have been written into a PROM chip, they are permanent and cannot be
altered. There are other types of ROM control chips available, however, that can be
erased and reprogramed. Since one type of erasable and programable read-only memory
(EPROM) chip needs to be removed from the processor and exposed for some time to
ultraviolet light before it can accept new contents, it is hardly suitable for use by
application programers. Another type of electrically erasable programable read-only
memory (EEPROM) chip is also available that can be reprogramed with special
electric pulses. Regardless of the type of ROM chip used, however, they all serve to
increase processor efficiency by controlling the performance
of a few specialized tasks.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
l-The first paragraph mainly discusses .......... .
a. the language of the computer
b. the language of the user
c. the use of a compiler
d. the use of an interpreter
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2. As we understand from the text, ..........
a. the user does not have access to the arithmetic/logic memory
b. the user does not have access to the main memory
c. bipolar semiconductor chips cannot be used in the primary storage
d. metal-oxide semiconductor chips cannot be used in the ALU
3. The third paragraph mainly describes ........... .
a. random-access memory chips used to store and retrieve data and
instructions when necessary
b. semiconductor storage elements including the types used in the
primary storage and those used in other sections of the processor
c. modern computers with semiconductor storage elements used in
their primary storage section
d. bipolar semiconductor chips used in the arithmetic/logic and other
sections of the computer
4. The fourth paragraph mainly describes .......... .
a. what dynamic RAM chips contain
b. how dynamic RAM chips are refreshed
c. the volatility of dynamic RAM chips
d. the characteristics of dynamic RAM chips
5. An uninterruptible power system .......... .
a. acts as a transistor being either on or off
b. acts as a capacitor storing either Os or 1s
c. is used as a backup to compensate for the volatility of RAM chips
d. is used as a backup to personal computers to raise their storage capacity
6. It is true that .......... .
a. the most developed processors have either specialized or general-purpose storage
b. the most basic computer functions are carried out by either highspeed
circuitry or by special regenerator circuits
c. computer users have no control over specialized storage elements
d. computer users are conscious of the operations performed by cache
7. As we understand from the text, .......... .
a. a microprogram does not efficiently perform high-level operations
b. a microprogram is a sequence of elementary instructions residing in
the processor
c. ROM chips are interchangeable with static RAMs
d. EPROM chips cannot be reprogramed by application programers
B. Write the answers to the following questions.
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1. What is the function of a compiler?
2. What is a word ?
3. Who decides on the number of bits per word?
4. What does RAM stand for?
5. What is the major advantage of random-access memory chips?
6. What are the similarities and differences between the dynamic and static RAM
7. What is the function of cache memory?
8. What are microprograms used for?
9. How do ROM chips differ from RAM chips?
10. How does the computer user customize a system?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Organizing Data for Processing
Data is the raw material to be processed by a computer. Such material can be letters,
numbers, or facts-such as grades in a class, baseball batting averages, or light and dark
areas in a photograph. Processed data becomes informa tion-data that is organized,
meaningful, and useful.
To be processed by the computer, raw data must be organized into characters,
fields, records, files, and data bases.
A character is a letter, number, or special character (such as $, ?, or *). One or more
characters comprise a field.
A field contains an item of data. For example, suppose a sports club was making
address labels for direct mailing. For each person, it might have a date-joined field, a
name field, a street address field, a city field, a state field, and a postal code field.
A record is a collection of related fields. Thus, on the sports club list, one person's
date-joined, name, address, city, state, and postal code would comprise a record.
A file is a collection of related records. The entire list of address labels for the
sports club would be a file.
A data base is a collection of interrelated data stored together with
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Minimum redundancy. Specific data items can be retrieved for various
applications. For instance, the sports club data could be obtained according to
state or postal code or alphabetically by last name.
Whenever a change is to be made to stored data, a record is generated
containing the new data. The record is called a transaction. Whenever files
are changed to reflect new information, the process is called updating. Files of
records are stored on some form of medium, usually magnetic disk or
magnetic tape, so they can be read into main computer storage for processing.
A file can be a transaction file, one that contains modifications to
existing records. For example, in our address label list, a transaction would be
a change in a label (a new address), an added label (a new member), or a
deleted label (a member resigns). Or a file can be a master file, which
contains relatively permanent data-the master address label list, in this
case-that is updated by a transaction file.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1. cache
2. compiler
3. customize
4. encounter
6. execute
7. fetch
8. interpreter
9. leak
10. mainframe
11. manipulate
12. metal-oxide semiconductor
13. microprogram
14. MOS
15. PROM
16. random access
17. redundance
18. static RAM
19. transaction
20. ultraviolet
21. vendor
22. volatile storage
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Section One: Reading Comprehension
Data Processing
The arithmetic and logic unit (ALU) of a computer may be thought of as the
heart of a computer since this is the component which performs the data
manipulation operations which are essential in any data processing task.
Similarly the control unit may be likened to the brain of the computer since its
function is to control the program execution according to the sequence of
instructions encoded as the computer program. The control unit does this by
[etching the instructions one by one from the memory and then obeying the
data manipulation which each instruction specifies. Thus each instruction
execution may be seen to require a two-beat cycle, where the first beat always
fetches the instruction and the second beat executes the function specified by
the instruction. The fetch operation is invariant for all instructions, but the
execute operation varies according to the instruction fetched and may, for
example, cause the addition of the contents of two registers, the clearing of a
memory location or the transfer of data between a register and a memory
location. The computer central processing unit or microprocessor unit must
therefore perform three tasks, each of which is examined in more detail
CPU Communication With Memory
Data are communicated between the CPU and the memory of a computer
using buses; similarly, communication between the registers in a computer
CPU occurs via buses internal to the CPU. Conceptually there is no difference
between these two types of bus, but in microprocessors the external buses
connecting to memory are brought out to the pins of the microprocessor while
the internal buses exist only on the silicon chip, and are inaccessible
In order for the CPU to be able to access instructions or data stored in the
main memory, the CPU must first of all supply the address of the required
memory word using an address bus. When this address has been specified, the
required instruction or data can be read into the CPU using the data bus.
In the case of data it may also be necessary to write data into the
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memory from time to time as well as to read data from the memory. One way
to do this would be to have two separate data buses, one for reading data from
the memory into the CPU, and one to write data back into memory. To
minimize the number of pins required on a microprocessor CPU chip ,
however, it is common to multiplex the data read and data write buses into a
single bus with a separate read/write control line which specifies the direction
of data transfer.
A further requirement of the CPU is a register to keep track of which
instruction is being executed, so that instructions can be executed in sequence.
This register is called the program counter (also sometimes known as the
instruction counter), and acts as a pointer to the next instruction to be
executed. It is incremented immediately after each instruction is accessed
from memory, and also after accessing every instruction operand or operand
address if the operand is stored in a separate word of memory. Each time that
it is required to read the next instruction word from memory, the output of the
program counter is connected via the address bus to the memory, thus
supplying the address of the next instruction. After allowing an appropriate
time for propagation delays the instruction word can be read into the CPU
from the data bus.
Another register called the accumulator is also provided for short-term
data storage in the CPU.
Instruction Execution-Data Manipulation
If it were only possible to use a computer to transfer data backwards and
forwards between registers and memory, the computer would be of very little
practical use; its power lies in its ability to manipulate and process
information so that new information can be generated from data which already
exists. The arithmetic and logic unit (ALU) is the component of a computer
system which performs the task of data manipulation. Generally it is adequate
to consider the ALU as a black box with two one-word-wide data inputs, and
one one-word-wide data output. In addition, the ALU contains a number of
control inputs which specify the data manipulation function to be performed.
The ALU is a combinational logic circuit, whose output is an instantaneous
function of its data and control inputs; it has no storage capability. Thus the
result of any ALU operation must be stored in an accumulator.
Instruction Interpretation and Control
A final aspect of the CPU to be considered in this unit is the mechanism for
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converting the instruction opcodes into the low-level control signals which
cause the operation specified by the opcode to be executed. This is the
function of the control unit of the computer. Because the instruction sets of
different types of computer differ, so, too, do the control units differ in the
way that they decode the operation codes to elemental control signals for the
rest of the computer. Nevertheless, some observations may be made about the
general structure of computer control units.
When the instruction opcode is loaded into the CPU from the computer's
memory, it is stored in the instruction register. This register provides
temporary storage for the opcode while it is applied to the inputs of an
instruction decoder, which converts the opcode to the required low level
control signals. The instruction decoder is a combinational logic component,
and may be implemented using random logic, a PLA, or even using a ROM,
though this tends to be rather inefficient in this application.
The combinational logic control signals are connected to all the
combinational components in the computer, such as the ALU function and
mode select inputs, and multiplexer and demultiplexer control inputs. The
timing of the signals applied to these units is unimportant, so long as sufficient
time is allowed for data and control signals to propagate through all
The register control outputs of the instruction decoder are connected to
the write (latch data) inputs of all the CPU registers, and to the memory
read/write line of the control bus. The timing of these signals is important
because it controls when data are written into the CPU registers or memory.
Part I. Comprehension Exercises
A. Put "T" for true and "F" for false statements. Justify your
…….1. The control unit directs and coordinates the computer system in
executing stored program instructions.
…….2. The control unit applies a separate fetch operation to each instruction.
…….3. The execute operation performed by the control unit is constant for all
…….4. The data read and the data write buses are multiplexed in order to
reduce the number of pins required on a microprocessor.
…….5. As we understand from the text, data and information are the
…….6. The instruction decoder interprets the opcode received.
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B. Choose a, b, c, or d which best completes each item.
1. The first paragraph mainly ………….
a. introduces the control unit as the most active component of a
b. introduces the ALU and the control units as the most important
features of a computer
c. describes the mechanism of the ALU
d. describes the fetch and the execute operations
2. In order for the CPU to read data from or to write data into memory
a. the registers in the CPU must supply the required data
b. the registers in the CPU must be cleared of data
c. data and address buses must be used respectively
d. address and data buses must be used respectively
3. The fifth paragraph mainly discusses ………….
a. the mechanism of the program counter
b. the mechanism of the central processing unit
c. the order of instructions stored in memory
d. the order of instructions fetched from memory
4. Data manipulation is performed in ………..
a. the control unit
b. the arithmetic-logic unit
c. the main memory
d. the accumulator
5. It is true that ……..
a. the combinational logic control signals are connected to
certain combinational components in the computer
b. the combinational logic control signals are time dependent because they
control data written into the CPU registers
c. the instruction register provides temporary storage tor the operation part of
an instruction received from memory
d. the instruction decoder receives the operation part of the instruction
interpreted in the instruction register
C. Answer the following questions orally.
1. What is the function of the ALU?
2. Why is a two-beat cycle needed for each instruction execution?
3. How does the CPO communicate with memory?
4. What are the data and the address buses'!
5. What is the function of a read/write control line'!
6. What is a program counter?
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Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. An………. specifies an operation and the values or locations of its operands.
a. instruction
b. integrator
c. instruction register
d. instruction code
2. The ………. receives the results of arithmetic manipulations within
the CPU.
a. data bus
b. address bus
c. accumulator
d. multiplexer
3. The microprocessor……… connect to other system components such as
memory and input/output interfaces, therefore their characteristics must be
carefully defined.
a. accumulators
b. program counters
c. internal buses
d. external buses
4. Most CPUs contain a number of which provide short-term storage for the
results of arithmetic and logic processing.
a. instruction counters
b. stack pointers
c. general-purpose registers
d. special-purpose registers
5. The ……… converts the machine code instructions fetched from the
computer's main memory into the actual control signals required to execute
the function specified by the instruction.
a. arithmetic-logic unit
b. control unit
c. program counter
d. instruction register
B. Fill in the blanks with the appropriate form of the words given.
1. Compute
a. A ………… is an electronic symbol-manipulating system designed and
organized to automatically accept and store input data, process them, and
produce output results under the direction of a detailed step by step stored
program of instructions.
b. Addition, subtraction, multiplication, and division are that the computer can
perform. ,
c. The computer component that performs the mathematical opera- tions
required for problem solving is called the element.
2. Address.
a. Symbolic addressing is the practice of expressing an ……., not in
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terms. of its absolute numeric location, bur rather in terms of symbols
convenient to the programmer.
b. The computer memory contains ………… locations.
c. A 16-line bus is built into a microprocessor chip to determine the primary
storage locations of the needed instructions and data.
d. To improve the data handling and capabilities of their products,
early 1980s. !
3 .Differ
a. There is not a very big …………. in flowcharting for a program to be
written in Cobol or Fortran.
b. There are many ………..computers manufactured today, and a buyer must
be able to compare the advantages and disadvantages of each.
c. The opinions of programmers as to the best way of solving a problem
often …………... greatly.
4. Logic
a. A program must be…………. organized if successful results are to be
b. There are three basic kinds of …………….operations: equal to, less than,
and greater than.
5. Store
a. In a personal computer system, operating-system programs are
commonly……………… on a floppy disk.
b. Generally, the larger the system, the greater its …capacity is.
c. Small magnetic tapes and floppy disks are typically used for offline secondary ………….
d. One of the basic types of processing a computer can do is the
…..... of data.
C. Fill in the blanks with the following words.
Although microcomputers have the general characteristics of digital
………, a notable property of microcomputers is their relatively ……….cost
and small size. This has greatly …….to their popularity and success. While
large ………computers and minicomputers have more computational power
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…………. , this power is not always needed in every ………… .
Furthermore, the cost of large computers and…………. has often precluded
their adoption into certain……… systems which would have otherwise
benefited from their ………..The microprocessor has now made it possible to
use a/an ………device in a logic system where cost constraints ……...than
speed and computational power are …………. .
D. Put the following sentences in the right order to form a paragraph.
Write the corresponding letters in the boxes provided.
a. This was a major advance in the electronics industry, but it only paved the
way for greater things: the evolution of the many analog circuits and the more
complex digital circuits which contained several transistors and other
b. The number of components which can be stacked on one chip determines
the complexity of the circuit.
c. Because these MSI. circuits are more reliable, cost less, and use .less power
than the transistor equivalents, they are easier and cheaper to use.
d. The microprocessor evolved from the transistor, which was the first major
semiconductor device.
e. Techniques were developed to put more and more complex circuits on the
same chip, evolving into medium scale integration (MSI).
f. Several transistors were then put on one semiconductor substrate, and the
integrated circuit (IC) evolved.
Section Two: Further ReadinK
System Software
A vital part of any general-purpose computer is the system software, or
software tools which are used in conjunction with the computer hardware.
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Without the system software the computer is rather like a car without petrol;
although the basic mechanics of the system exist, there is no way of actually
using it. This unit concentrates on the software tools which are required to
turn a general-purpose computer into a useful computer system for applications
programming, for microprocessor applications development, or for use as a
business or administrative system. Some insight into the types of programs
required can be gained by considering the various tasks which might be
undertaken using a general-purpose computer.
The first requirement when a computer is switched on is some kind of loader
program which can be used to load any other program from the backing
storage medium into memory prior to execution. In most modern computers
the loader would be a rather large program stored in ROM and designed to
read programs from a disk, as shown in Figure 8-1.
The concept of using one piece of software to enable another (generally
more complex) piece of software to be run is called bootstrapping. Thus the
loader program provides a method of bootstrapping other more complex
programs. This approach is taken to minimize the use of the computer system
memory; the loader program which resides permanently in ROM typically
occupies only about 1 kbyte of memory, and therefore frees most of the
Top of memory
Bottom of memory
Figure 8-1. Operation of a Loader Program.
remaining memory for other programs. The loader program is also often
combined with a simple monitor or debugger program which may be used to
debug machine code programs, and also to verify operation of the computer
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hardware without requiring access to a backing storage medium.
A different type of computer may also be used for bootstrapping, by making
use of cross-software. For example, a program may be written for a
microprocessor such as MC6809 on a minicomputer or mainframe computer,
using an assembler written in a language which is available on the large
computer, such as PASCAL, C, or FORTRAN. The assembler program which
runs on the minior mainframe computer is known as a cross-assembler
because it produces machine code for the microcomputer and not for the
computer on which the assembler is run. The machine code program can then
be transferred to the target computer either in PROMs or via any of the other
storage media, or it can be downloaded directly via a serial or parallel
communication link.
Disk Operating System
Software which enables programs to be loaded from (and stored on) backing
storage media can be combined with facilities which handle the display
terminal(s) and other computer peripherals such as printers, plotters, and so on,
to provide a general-purpose control program. This operating system program
generally makes use of floppy or hard disks as the main backing storage
medium, because of the ease of randomly accessing different areas of disk
containing different files. Hence the program is known as a disk operating system
A disk operating system provides two fundamental facilities. First, it
provides a mechanism for communication between the computer and the user by
handling input and output to the user's console and executing the commands
specified by the user. Second, it provides a mechanism for program storage and
retrieval, though generally in a rather more flexible and sophisticated form than
the simple loader program. Through the medium of the operating system users
are able to access source files and call up text editors to make changes to
them. They can then assemble or compile the source code (as appropriate) to
produce machine code programs, and the machine code can be loaded into
memory and executed, all under the control of the operating system.
Files, which may contain machine code, ASCII coded source text, data or any
other information, are stored on the disk in a format defined by the operating
system, and are accessed via filenames. The filenames, which are simply
mnemonics chosen by the user to reflect the contents of a file, are stored in a
directory which can be examined by the user. Manipulation of files
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can then be performed simply by referring to the file using its filename. In a
similar way, all the other system programs described in this unit are accessed
through an operating system.
The operating system also provides a number of housekeeping utility
programs which are necessary so that the user can maintain his files on the
disk in an orderly fashion. Typically, the operating system includes the
(a) Directory listing so that the user can determine what files are on the
disk, their size, when they were created and other useful information.
(b) File erasing so that unwanted files may be removed from the disk to
free space for other purposes.
(c) File renaming so that filenames may be changed if required.
(d) File transfer so that files may be copied to another disk for backup
purposes or for duplication.
(e) File listing so that the contents of text files may be printed.
(f) File execution so that machine code files may be loaded into memory
and executed.
In addition, operating systems may include many other more advanced and
specialized facilities which are beyond the scope of this discussion.
Finally, the operating system provides a straightforward mechanism to
enable the user to access terminals, printers and storage devices connected to
the computer from within his application programs. Access is achieved by
executing subroutine calls to standard subroutines within the operating system
which control these devices. This saves the user from needing a detailed
knowledge of the interface characteristics of the peripheral devices connected
to the computer if he simply wishes to run application programs which access
the standard computer peripherals under the control of the operating system.
Text Editor
Given a loader for bootstrapping the operating system when the computer is
switched on, and the ability to manipulate files stored on a disk storage
system, the user's next requirement is generally a facility for developing
application programs in a high-level language or assembly language. In either
case, the application program is initially written as a source code text file
using a text editor. The editor program must therefore provide a mechanism
for inputing the source code program and saving it on disk, and subsequently
facilities must also be available for making changes, corrections and additions
to the source program.
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Three types of activity are involved in using a text editor. First, the user
must identify the text which is to be manipulated, or the place at which text is
to be inserted in the file. This requires a pointing operation to specify where
the change is to be made. Second, the user must specify the operation or
command which is to be performed. Typical editing operations will include
insertion of new text, deletion of text, and substitution of new text for old. In
each case, the change may be made on a character, word, or line basis. More
sophisticated editors also include commands to search a complete file for
occurrences of a specified character string, and optionally replace this with a
new string. They may also provide facilities for merging separate files or
subdividing files. Finally, having specified the required editor command, the
user must type in the replacement text or new text if the command requires it.
Computer text editors can be categorized into two broad types: lineba3'ed and screen-based. Screen-based editors are the simplest to use because
the display terminal represents a window on the text file, and the pointing
operation is achieved by using cursor keys on the terminal to position the
cursor where a change is to be made. If the user wishes to move outside the
range of the current display window, the screen is scrolled up or down to
position the window at the required place in the file. Commands are generally
displayed using a simple menu positioned on a fixed 'status' area of the screen,
and selected by typing the first letter of the command or using special
command function keys. Thus the text editor is simple and straightforward to
use, and the effect of executing a command is immediately obvious on the
If a screen-based editor cannot be used, the alternative is to use a linebased editor. All operations in this case are referred to a specified line in the
source file. The line may be specified using a line nulJlber, or it may be
identified using an invisible cursor. In this latter case, commands are executed
to move the invisible cursor through the file, and the current cursor position
within the file can be determined by causing the editor to list the current line
on the terminal. Line-based editors generally require much greater memory on
the part of the user and hence require significant training before competence is
Assemblers and Compilers
Once a source file has been prepared using a text editor, the next stage is to
convert the program to machine code using an assembler (if the program is
written in assembly language) or a compiler (if in a high-level language). For
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short programs, the facilities introduced so far prove quite adequate; as larger
and more complex programs are written, however, some disadvantages begin
10 become apparent.
As the size of the source program increases, the time required to assemble
or compile the program also increases as does the editing time (since it takes
longer to locate the required part of the program within the file). Thus the
efficiency of the programmer begins to fall. One way to resolve this problem
is to split the program into a number of separate source modules, each of
which can be edited and assembled or compiled separately. A mechanism is
then required to build the various program segments together into a final
machine code program. Another related problem also may become apparent;
this is a need to be able to segment and separate different parts of a program
physically in memory. As an example, in developing a microprocessor
application, it may be required to specify one memory area for program and
data constants which will be stored in ROM, and a separate non-contiguous
area for data variables stored in RAM. Furthermore, at the time the program is
written, the final allocation of memory may not be known. Thus a simple
mechanism for relocating the program to different memory addresses is
Debuggers and Simulators
Once a program has been written and assembled or compiled, the next stage is
to run the program and verify its operation. Even programs written by
experienced and expert programmers very seldom work correctly to begin
with. It is very difficult to foresee all the ways in which a program algorithm
is executed and invariably algorithm faults are shown up when the program is
first tested. To cater for this, debuggers and simulators are required to test a
If the program is written in the computer's native assembly language, then
it is often possible to load the machine code program into the computer's
memory and execute it under the control of a monitor program of the type
already discussed. If, however, the program has been written using a crossassembler, then it cannot be executed directly, but instead a simulator
software package is often provided which interprets each target machine code
instruction and modifies the memory locations in the host computer to
simulate the action of the program on the target computer registers, memory
and input/output circuits. If the configuration of the application computer is
very different from that of development computer, it may not be possible fully
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to debug the software within the development system. In this case, in-circuit
emulation facilities are used to complete the debugging process.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. A loader program is designed ………….
a. to read the programs executed in the main memory into the secondary
b. to read the program to be executed from the secondary storage into the
main memory
c. to bootstrap pieces of simpler programs in order to back up the
computer memory
d. to bootstrap more complex programs in order to verify their validity
2. A program written for a microprocessor may be assembled on a mini or
mainframe and then transferred to the target microprocessor by ……….
a. a loader
b. a monitor
c. a text editor
d. a cross-assembler
3. As we understand from the text,………
a. operating system programs are commonly stored on disks
b. tapes are easier to use as the backing storage medium
c. a disk operating system is not of much help to the user
d. a disk operating system is not as efficient as a loader
4. An operating system enables the users ……….
a. to communicate with the computer
b. to make changes to the source files
c. to have access to all the system programs
d. all of the above
5. It is true that .......... .
a. the user can refer to files via filenames
b. the user cannot manipulate the source files
c. the operating system makes the system more difficult to use
d. the operating system does not control the overall operations of the
6. The text editor …………….
a. controls the execution of other programs
b. allows a user to enter and store programs
c. translates each source language statement into a sequence of machine
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d. inputs the source code programs, saves it on disk, and reviews and alters
its materials
7. Using a text editor, the user must ……….
a. identify the text to be manipulated
b. specify the operation to be performed
c. type in the new text if required
d. all of the above
8. According to the text, ……….
a. a window is a portion of the visual display screen used to show the
current status of an application of interest
b. a cursor is positioned where a change is to be made when a line- based
editor is used
c. line-based editors require smaller memory compared with screen- based
d. screen-based editors are hard to use since the user cannot move outside
the range of the current display window
9. It may be inferred from the text that ………..
a. screen-based editing is not necessarily carried out under the user's
b. screen-based editing requires a high communication rate between the
console and the computer
c. the concept of a cursor seems meaningless on a display terminal
d. the concept of a cursor is only meaningful when a change is made on a
10. An application program, once prepared, has to be ………
a. converted into a high-level language
b. segmented physically in memory
c. split into a number of modules
d. assembled or compiled .
1. What is bootstrapping?
2. Why does bootstrapping minimize the use of the computer system
3. What is a directory?
4. What does an operating system include?
5. What is a filename?
6. Why does the user not need to know enough about the interface
characteristics of the peripheral devices connected to the computer?
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7. What is a pointing operation used for?
8. What do editing operations include?
9. What problems arise as the size of the source program increases?
10. What are debuggers and simulators used for?
Section Three: Translation Activities
A. Translate the following passage into Persian.
The Microprocessor
One of the results of the advancement in solid-state technology is the
capability of fabricating very large numbers of transistors (say, 1000 and over)
within a single silicon chip. This is known as large-scale integration. A direct
consequence of large-scale integration is the microprocessor. In general, a
microprocessor is a programmable logic device fabricated according to the
concept of large-scale integration. A microprocessor has a large degree of
flexibility built into it. By itself it cannot perform a given task, but must be
programmed and connected to a set of additional system devices. These
additional system devices usually include memory elements and input/output
devices. In general, a set of system devices, including the microprocessor,
memory, and input/output elements, interconnected for the purpose of
performing some well-defined function, is known as a microcomputer or
microprocessor system.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1. administrate
2. bootstrap
3. console
4. cross-assembler
5. cross-software
6. debugger
7. file erasing
8. invisible cursor
9. line-based editors
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Section One: Reading Comprehension
Computerized Critical Care Areas
critical care unit is an area in a hospital where highly trained personnel
and sophisticated equipment are concentrated to take care of a limited
number of actually or potentially severely ill patients. Special units may be
known as general intensive care units or labeled according to the type of
patients treated. Thus we have medical, surgical, neurological, respiratory, or
pediatric intensive care units. Coronary care units, recovery rooms, telemetry
monitoring areas, burn and trauma units are also critical care areas. Physicians
trained to work in these units are commonly referred to as ‘intensivists’.
Likewise, nursing personnel permanently assigned to these units undergo
specialized training.
Computers are now commonplace in critical care areas in both large and
small hospitals. They cover a wide range of applications, from the microprocessor that controls specialized bedside and nurses’ desk monitoring
equipment to the mini- or mainframe computer that is part of either a
dedicated critical care system or an integrated overall hospital-wide information facility. Figure 10-1 represents a commonly seen special care unit
arrangement, consisting of (a) microprocessor-controlled patient monitoring
hardware with bedside and nursing-desk scopes and controls, as well as
hard-copy functions, and (b) remote computer stations, usually video display
terminals (VDT) communicating through a central mainframe installation
with ancillary departments, other patient care areas, and business and
financial offices.
The first attempts at computer-assisted patient monitoring in critical
care areas took place in the 1960s. Some of the early applications were based
on electrocardiographic waveform analysis and attempted to establish the
morphologic diagnosis of myocardial ischemia or injury, conduction defects,
or chamber enlargement. Other developments focusing on the automated
recognition of cardiac arrhythmias followed but had limited success, and to
this day arrhythmia interpretation by computer remains an elusive goal since
technology has not yet equalled the human mind in recognizing complex
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Figure 10-1. Diagram of a Typical Critical Care Unit Arrangement Consisting of
Microprocessor-Controlled Bedside and Central-Station Monitors, Complemented by a
Computer Station.
Today it is unusual to find drug infusion devices, EGG and blood
pressure monitors, intraaortic balloon assist pumps, or other critical care unit
devices that are not controlled by microprocessors. Many of these new devices
also have built-in communication controllers that allow them to transfer
information to, and/or be controlled by, an external computer system. A
microprocessor-controlled bedside physiologic monitoring unit may be used in
a coronary care unit. It is primarily used to acquire, display, and transmit a
patient's heart rate, electrocardiogram, and arterial blood pressure, but
additional parameters can be incorporated. Built-in audible and visual alarms
alert the staff if preset upper or lower limits are exceeded in any monitoring
Channel .
With the advancement of technology, personal computers and even small
hand-held computers have become popular because they offer powerful
processing tools in small and relatively inexpensive packages. Programs are
available that accept as input homodynamic and blood gas information to
calculate and print a patient's homodynamic profile. Other programs control
infusion of drugs or blood, perform cardiac output and drug dosage calculations, help manage the treatment of acid-base disorders, or assist in hyperalimentation therapy. These small dedicated computers allow, in. some cases,
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transfer of information to a larger central system if this is available.
In addition to automated bedside monitoring, computer technology is
used in special care units for acquisition, storage and processing of patient
information needed for determination and reporting of trends, requesting and
reporting ancillary department tests or procedures, and administrative functions such as inventory control and billing. These uses result in the elimination of manual tasks, improve accuracy, eliminate redundant bookkeeping,
and reduce the possibility of human errors.
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
……. 1. A seriously ill child is usually treated in a pediatric intensive care
……. 2. Critical care areas in hospitals are rarely computerized.
……. 3. Computers have been used in critical care areas since they first
appeared on the market.
……. 4. Arrhythmia interpretation has successfully been performed by the
computer since 1960.
……. 5. Small hand-held computers may also be used in critical care units.
……. 6. Small computers in critical care units are used for a wide variety
of applications.
B. Choose a, b, c, or d which best completes each item.
1. The first paragraph mainly describes ………. .
a. physicians who work in critical care units
b. patients who are treated in critical care units
c. the critical care areas in a hospital
d. the sophisticated equipment used in a hospital
2. Compared to human beings, computers ………. .
a. are superior in mental ability
b. are inferior in mental ability
c. can better interpret complicated heart patterns
d. can better recognize heart arrhythmias
3. As we understand from the text, critical care units have been established in the hospitals………… .
a. for more than 35 years
b. for less than 35 years
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c. since the computers could analyze waveforms
d. since the computers were produced
4. Paragraph three mainly discusses………. .
a. electrocardiograph waveform analysis by the computer
b. diagnosis of myocardial injury and conduction defects by the computer
c. some applications of the computer in critical care areas
d. some applications of the computer in the recognition of cardiac
5. It is true that ……… .
a. computer-controlled devices are more effective than intensivists in the
recognition of cardiac arrhythmias
b. computer-controlled devices have not been as effective as expected in
critical care units
c. many of the new care unit devices are controlled by the nursing
d. most of critical care unit devices are now computerized
C. Answer the following questions orally.
1. What is a critical care unit?
2. What are some examples of intensive care units?
3. What is the function of built-in communication controllers in a
computerized critical care unit?
4. How is a built-in audible alarm helpful?
5. What are some of the applications of small computers used in critical
care units?
6. What does the last paragraph mainly discuss?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. Patients in critical care units are taken care of by……..
a. surgeons
b. ophthalmologists
c. intensivists
d. pediatricians
2. The objective of is to register graphically movements of the heart.
a. a cardiographer
b. a cardiograph
c. a transducer
d. a simulator
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D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. Less critical files should be backed up at intervals dictated by the
frequency of updates to those files.
b. Thus a fire, flood, or other natural disasters cannot cause total loss of
c. The data contained in a medical database are vital to the institution or
the medical office, and their safety and integrity must be preserved.
d. In addition, a standard operating procedure should be instituted whereby
weekly or monthly backup files are created and stored in removable
storage media, kept preferably at a location remote from the institution
or office.
e. Critical data files, shared by many users and updated frequently, should be
backed up on removable storage media at regular and frequent
Section Two: Further Reading
Planning and Designing a Computerized
Critical Care Unit
Experience has shown that a computer system can help reduce the length of
stay of a patient in a special care area. In this age of increased awareness of
the cost of providing high technology health care, this fact results in better
utilization of the resources in the unit, and lower costs can be expected. Those
individuals involved in the process of planning, designing, and activating a
state-of-the-art critical care unit incorporating advances in computer technology may want to follow the steps outlined below. However, every institution
presents a different environment, and, therefore, individual designs will
probably be significantly personalized.
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1. A planning committee including medical and nursing staff members as
well as high-level administrative personnel should be established. Delegating
the task of planning a project of this magnitude to lower-level management
may not produce satisfactory results.
2. A comprehensive evaluation of existing and projected patient care
needs should be undertaken to determine specifications for the special care
unit in question. Some of these specifications may be determined by existing
facilities and/or budgetary constraints. A state-of-the-art unit should incorporate the capability to acquire and process signals representing biological
variables as continuous functions of time (physiologic monitoring), as well as
acquisition, storage, processing, and recall of discrete patient information. If
no centralized computer system is available in the institution, then a small
dedicated special care system may be a realistic approach. On the other hand,
if a comprehensive hospital information system is currently available, then a
major objective should be the integration of the special care unit into the
central system.
3. When considering automation in the unit, a decision must be made
early in the planning process whether to obtain a commercial turnkey
hardware-software package, as opposed to the in-house development of a
dedicated microprocessor or minicomputer-based special care unit system. If a
mainframe central installation is available, a link to it should be considered in
cither case. Figure 10-2 lists several approaches to linking a mainframe-based
hospital information system and a dedicated local or satellite computer.
The choice of one of these approaches may depend to a large degree on
existing facilities, equipment, staff, and experience. In-house developed
computer systems have the advantages of being designed and built according
to the needs and desires of the staff and afford the ability to make changes as
they become necessary, sometimes on short notice. One disadvantage of the
in-house approach is that the development time may be long, and, therefore,
personnel costs may be high. Commercially available turnkey systems, on the
other hand, are usually ready for production work once installation is
completed, and development time and costs are substantially lower. However,
modifications to the system to meet existing institutional policies or
procedures, if needed, may be expensive or not possible. This rigidity of design
entails, in many cases, modifications in policies or procedures to conform to a
somewhat inflexible, commercial package.
These decisions may not be simple but should be based on the approach
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that best suits the present and future needs in the existing hospital
environment. In any case, provisions should be included during the planning
stages for future implementation or expansion of capabilities for automated
entry, communication, archiving, processing, and reporting of information.
Shared memory
Telephone modem
Local area network
Off-line magnetic tape
data transfer
Off-line floppy disk
data transfer
.. Off-line removable
hard disk data |
Microcomputer interface
Remote Station
RS232 hardware interface
Figure 10-2. Direct (a) and Indirect (b) Approaches to Establishing
Communications Between a Dedicated Special Care System and a Central
Mainframe-Based Hospital Information System.
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4. Whether acquiring a commercial package or designing an in-house
system, the administrative aspects of the operation of the special care units
should not be neglected. A special care unit package should provide
administrative services (or interact with any existing system that already
provides them). These should include, but not be limited to, inventory control,
patient charges, bed use, and cost analysis information with daily, weekly,
and/or yearly reports and appropriate audit trails.
5. Once the new computerized techniques for data acquisition, storage,
processing, and reporting become established, usually alter a suitable
‘paralleling’ period, the old manual methods should be discontinued. How ever, contingency procedures based on the old methods should be established,
documented, and tested frequently in the event of an equipment breakdown or
other computer system failures.
6. If not already available, an uninterruptible power supply (UPS) should
be included as an integral part of the state-of-the-art special care unit. A UPS
system should provide power to all computer systems in the unit, if the
normal ac service is interrupted. Many computer storage devices (i.e.
random-access memory chips) are volatile and do not retain information if the
power is interrupted, even for a fraction of a second. Therefore, there is a real
potential for losing critical patient information obtained during spontaneous
clinical events that cannot be reproduced. UPS systems are available in many
configurations and capacities depending on the particular electrical service
required. Typically, they include a utility-fed rectifier that supplies dc power to
a set of batteries and an inverter that provides clean, transient-free ac power
to the equipment. The batteries provide backup during short power failures
(minutes) or until the hospital emergency generators take over in case of
longer outages.
The actual design of the state-of-the-art critical care unit follows the
planning stages and should be a multidisciplinary task. Physicians, nurses,
architects, clinical and/or biomedical engineers, data processing personnel, and
systems engineers should integrate the design team.
Often, too little thought is given to the practical aspects of the room
layout, including computer cabling and connections and the design and
placement of computer terminal cabinets. These items are often ignored until
after the room is already under construction or completed. Many potential
problems can be eliminated by building an actual full-size prototype of the
proposed critical care area to ensure optimum placement and accessibility of
all monitors, computer-related equipment, and other necessary devices.
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B. Write the answers to the following questions.
1. Who should be involved in a planning committee?
2. How do constraints affect the plan specifications?
3. What is the difference between a commercial turnkey system and an inhouse developed computer system?
4. What are the advantages and the disadvantages of an in-house computer system?
5. Why is designing a critical care unit considered a multidisciplinary
Section Three: Translation Activities
A. Translate the following passage into Persian.
Selection of Monitoring Equipment
There are always questions regarding the number and types of biological
variables that should be continuously monitored on special care unit patients.
In most institutions, the choice is dictated largely by the capabilities and
limitations of the monitoring systems commercially available at any given time
It is desirable for potential users to become familiar with the technical
terminology on equipment specification sheets. These specifications usually
describe the actual capabilities of the equipment much more clearly and in
more detail than do aggressive sales persons or colorful, eye-catching sales
literature. If an in-house biomedical engineering department is available, it
should evaluate this information and help the medical and nursing staffs
interpret it.
Determining what systems are available may require a comprehensive
review of the scientific and trade literatures as well as calls and site visits to
vendors and users for detailed information. It is recommended that a firm
understanding of the capabilities and limitations of any system be established
before pricing and contracts are considered. It should be stressed that the
selection process should include site visits to institutions that have used or are
using equipment similar to that being considered. This should include visits to
institutions not in the vendors’ reference list.
Finally, it should be noted that the successful implementation of critical
care unit systems rests not only on the adequacy of the hardware and software
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but also on the human component. Users, including physicians, nurses, and
technicians, should not only be involved in the selection process, but should
receive comprehensive hands-on training in the use of the equipment before
and during the actual implementation phase.
B. Find the Persian equivalents of the following terms
expressions and write them in the spaces provided.
1. ancillary
2. arterial blood pressure
3. biomeical engineer
4. cardia
5. cardiac arrhythmia
6. coronary care unit
7. critical care unit
8. drug infusion device
9. homodynamic profile
10. hyperalimentation therapy
11. intensive care unit
12. intensivist
13. intraaortic balloon assist Pump
14. myocardial ischemia
15. neurologicall
16. pediatric intensive care unit
17. recovery room
18. respiratory
19. telemetry monitoring area
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Section One: Reading Comprehension
Introduction to Control Systems
Automatic control has played a vital role in the advance of engineering and
science. In addition to its extreme importance in space-vehicle systems,
missile-guidance systems, aircraft-autopiloting systems, robotic systems, and
the like, automatic control has become an important and integral part of
modern manufacturing and industrial processes. For example, automatic
control is essential in the numerical control of machine tools in the
manufacturing industries. It is also essential in such industrial operations as
controlling pressure, temperature, humidity, viscosity, and flow in process
In studying control engineering, we need to define those terms that are
necessary to describe control systems, such as plants, disturbances, feedback
control, and feedback control systems.
Plants. A plant is a piece of equipment, perhaps just a set of machine
parts functioning together, the purpose of which is to perform a particular
operation. In control systems, any physical object to be controlled such as a
heating furnace, or a spacecraft is called a plant.
Disturbances. A disturbance is a signal that tends to adversely affect the
value of the output of a system. If a disturbance is generated within the
system, it is called internal, while an external disturbance is generated outside
the system and is an input.
Feedback Control. Feedback control refers to an operation that, in the
presence of disturbances, tends to reduce the difference between the output of
a system and some reference input and that does so on the basis of this
difference. Here only unpredictable disturbances are so specified, since
predictable or known disturbances can always be compensated for within the
Feedback Control Systems. A system that maintains a prescribed
relationship between the output and some reference input by comparing them
and using the difference as a means of control is called a feedback control
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system. An example would be a room-temperature control system. By
measuring the actual room temperature and comparing it with the reference
temperature (desired temperature), the thermostat turns the heating or
cooling equipment on or off in such a way as to ensure that the room
temperature remains at a comfortable level regardless of outside conditions.
Servo Systems. A servo system (or servomechanism) is a feedback
control system in which the output is some mechanical position, velocity, or
acceleration. Therefore, the terms servo system and position- (or velocity- or
acceleration-) control system are synonymous. Servo systems are extensively
used in modern industry. For example, the completely automatic operation of
machine tools, together with programmed instruction, may be accomplished
by the use of servo systems. It is noted that a control system, whose output
(such as the position of an aircraft in space in an automatic landing system) is
required to follow a prescribed path in space, is sometimes called a servo
system, also. Examples include the robot-hand control system,where the robot
hand must follow a prescribed path in space, and the aircraft automatic
landing system, where the aircraft must follow a prescribed path in space.
Automatic Regulating Systems. An automatic regulating system is a
feedback control system in which the reference input or the desired output is
either constant or slowly varying with time and in which the primary task is to
maintain the actual output at the desired value in the presence of disturbances. There are many examples of automatic regulating systems, some of
which are the Watt’s flyball governor, automatic regulation of voltage at an
electric power plant in the presence of a varying electrical power load, and
automatic control of the pressure and temperature of a chemical process.
Process Control Systems.
An automatic regulating system in which the output is a variable, such as
temperature, pressure, flow, liquid level, or pH, is called a process control
system. Process control is widely applied in industry. Programmed controls
such as the temperature control of heating furnaces in which the furnace
temperature is controlled according to a preset program are often used in such
Closed-Loop Control Systems. Feedback control systems are often
referred to as closed-loop control systems. In practice, the terms feedback
control and closed-loop control are used interchangeably. In a closed-loop
control system the actuating error signal, which is the difference between the
input signal and the feedback signal (which may be the output signal itself or a
function of the output signal and its derivatives), is fed to the controller so as
to reduce the error and bring the output of the system to a desired value.
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The term dosed-loop control always implies the use of feedback control
action in order to reduce system error.
Open-Loop Control Systems. Those systems in which the output has no
effect on the control action are called open-loop control systems. In other
words, in an open-loop control system the output is neither measured nor fed
back for comparison with the input. One practical example is a washing
machine. Soaking, washing, and rinsing in the washer operate on a time basis.
The machine does not measure the output signal, that is, the cleanliness of
the clothes.
Adaptive Control Systems. The dynamic characteristics of most control
systems are not constant for several reasons, such as the deterioration of
components as time elapses or the changes in parameters and environment.
Although the effects of small changes on the dynamic characteristics are
attenuated in a feedback control system, if changes in the system parameters
and environment are significant, a satisfactory system must have the ability of
adaptation. Adaptation implies the ability to self-adjust to self-modify in
accordance with unpredictable changes in conditions of environment or
structure. The control system having a candid ability of adaptation (that is, the
control system itself detects changes in the plant parameters and makes
necessary adjustments to the controller parameters in order to maintain an
optimal performance) is called the adaptive control system.
Learning Control Systems. Many apparently open-loop control systems
can be converted into closed-loop control systems if a human operator is
considered a controller, comparing the input and output and making the
corrective action based on the resulting difference or error.
If we attempt to analyze such human-operated dosed-loop control systems
we encounter the difficult problem of writing equations that describe
the behavior of a human being. One of the many complicating factors in this
case is the learning ability of the human operator. As the operator gains more
experience, he or she will become a better controller, and this must be taken
into consideration in analyzing such a system. Control systems having an
ability to learn are called learning control systems.
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
…… 1. Automatic control is an essential part of modern manufacturing and
industrial processes.
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…… 2. Any machine which is being controlled is referred to as a plant.
…… 3. The mechanism of a servo system is different from that of a
position control system.
…… 4. A control system whose output is required to follow a prescribed
condition may be called a servo system.
…… 5. The automatic controller functions more effectively in an openloop control system.
…… 6. An adaptive control system is capable of accommodating unpredictable environmental changes, whether these changes occur
within the system or external to it.
…… 7. An adaptive control system is designed to modify the control signal
as the system environment changes, so that performance is always
optimal whereas the human operator recognizes familiar inputs
and can use past learned experiences to react in an optimal
B. Choose a, b, c, or d which best completes each item.
1. A disturbance of a system.
a. has a positive effect on the output
b. has an unfavorable effect on the output
c. increases the efficiency
d. controls the efficiency
2. Feedback gives an automatic control system the ability…….. .
a. to deal with the unexpected disturbances in the plant behavior
b. to deal with the predictable disturbances in the plant behavior
c. to maintain a steady relationship between the output and some
reference input
d. to maintain the actual value of a disturbance constant
3. It is true that ……… .
a. the mechanism of the automatic regulating system is based on that
of the process control system
b. the mechanism of the process control system is different from that of
the automatic regulating system
c. an automatic regulating system compares the actual value of the
plant output with the desired value
d. an automatic regulating system maintains the actual value of the
plant output at the desired value in the presence of disturbances
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4. In a closed-loop control system, the controller ………..
a. regulates the internal disturbances of the plant and keeps them
under control
b. provides information about the actual plant output
c. reduces the difference between the input signal and the output signal
and brings the output of the system to a desired value
d. determines the value of the error and reduces its effect on the system
5. We may infer from the text that open-loop control systems ……. .
a. should be used for systems in which unpredictable disturbances occur
b. should be used for systems in which the inputs are know n in advance
c. are more complicated than closed-loop control systems
d. are more powerful than closed-loop control systems
C. Answer the following questions orally.
1. What part has automatic control played in the advancement of
2. What is a plant?
3. What are the internal and external disturbances?
4. What is the function of a feedback control system?
5. What is a servo system?
6. What is the mechanism of an open-loop control system based on?
7. What is an adaptive control system?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each Item.
1. The most fascinating developments in adaptive control systems lie in
the areas of pattern recognition and …….. systems.
a. learning
b. operating .
c. analyzing
d. reading
2. An automatic …….. compares the actual value of the plant output with
the desired value, determines the deviation, and produces a control
signal that will reduce the deviation to zero or to a small value.
a. amplifier
b. sensor
c. controller
d. transformer
3. A maintains the plant output constant at the desired value in the
presence of external disturbances.
a. capacitor
b. compensator
c. resistor
d. regulator
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4. In everyday life, occurs when we are aware of the consequences of
our actions.
a. adaptation
b. regulation
c. feedback
d. control
5. Control systems without feedback are called.
a. closed-loop
b. open-loop
c. adaptive
d. learning
B. Fill in the blanks with the appropriate form of the words
1. Heat
a. For high current levels, an external pass transistor may be required
with sinks to reduce the effective thermal resistance.
b. A heat coil is a protective device that grounds or opens a circuit, or
does both, by means of a mechanical element that is allowed to move
when the fusible substance that holds it in place is ………. above a
predetermined temperature by the current in the circuit.
c. A heater connector is designed to engage the male terminal pins of a
or cooling appliance.
d. A heater transformer supplies power for electron-tube filaments or
……… of indirectly heated cathodes.
2. Adapt
a. An ………. System is capable of accommodating unpredictable
environmental changes, whether these changes occur within the
system or external to it.
b. The vagueness surrounding most definitions and classifications of
adaptive systems is due to the large variety of mechanisms by which
……... may be achieved.
c. When high ……… is called for most present-day requirements will be
met by an identification-decision-modification system.
3. Accomplish
a. The dynamic characteristics of a plant must be measured and
identified continuously. This should be ……… without affecting the
normal operation of the system.
b. When tied in with learning approaches , pattern-recognition techniques will ……… adaptive-learning control.
c. A business system may consist of many groups. Feedback methods of
reporting the ………. of each group must be established in such a
system for proper operation.
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4. Cool
a. A heat exchanger or …….. is used in rotating machinery to transfer
heat between two fluids without direct contact between them.
b. The ………. of regulator elements refers to the method used for
removing heat generated in the regulating process.
c. In an electron device, a metallic part or fin extends the ..……area to
facilitate the dissipation of the heat generated in the device.
d. Air may be used as a……….. to remove heat from a machine.
5. Change
a. Modification refers to the …….. of control signals according to the
results of the identification and decision.
b. If parameters are…….. rapidly, a procedure known as alternate
biasing is employed.
c. Adaptive control systems are designed to modify the control signal as
the system environment ………. so that performance is always
d. Feedback allows us to cope with a……… environment by adjusting
our actions in the presence of unforeseen events.
C. Fill in the blanks with the following words.
An advantage of the closed-loop …… system is the fact that the use of
……. makes the system response relatively insensitive to……. disturbances
and internal variations in system ……… . It is thus possible to use relatively
……. and inexpensive components to obtain the accurate control of a/an …….
plant, whereas doing so is impossible in the open-loop ……… .
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. Some systems may have multiple inputs and multiple outputs.
b. A system may have one input and one output.
c. Such a system is called a single-input, single-output control system.
d. An example of such multiple-input, multiple-output systems is a
process control system that has two inputs (pressure input and
temperature input) and two outputs (pressure output and temperature
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e. An example is a position control system, where there is one command
input (desired position) and one controlled output (output position).
Section Two: Further Reading
Examples of Control Systems
Speed Control System. The basic principle of a Watt's speed governor for an
engine is illustrated in the schematic diagram of Figure 11-1. The amount of
fuel admitted to the engine is adjusted according to the difference between the
desired and the actual engine speeds.
The sequence of actions may be stated as follows: The speed governor is
adjusted such that, at the desired speed, no pressured oil will flow into either
side of the power cylinder. If the actual speed drops below the desired value
due to disturbance, then the decrease in the centrifugal force of the speed
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governor causes the control valve to move downward, supplying more fuel,
and the speed of the engine increases until the desired value is reached. On
the other hand, if the speed of the engine increases above the desired value,
then the increase in the centrifugal force of the governor causes the control
valve to move upward. This decreases the supply of fuel, and the speed of the
engine decreases until the desired value is reached.
Robot Control System. Industrial robots are frequently used in industry
to improve productivity. The robot can handle monotonous jobs as well as
complex jobs without errors in operation. The robot can work in an
environment intolerable to human operators. For example, it can work in
extreme temperatures or in a high- or low-pressure environment or under
water or in space. There are special robots for fire fighting, underwater
exploration, and space exploration, among many others.
The industrial robot must handle mechanical parts that have particular
shapes and weights. Hence, it must have at least an arm, a wrist, and a hand.
It must have sufficient power to perform the task and the capability for at
least limited mobility. In fact, some robots of today are able to move freely by
themselves in a limited space in a factory.
The industrial robot must have some sensory devices. In low-level robots,
microswitches are installed in the arms as sensory devices. The robot first
touches an object and then, through the micros witches, confirms the existence
of the object in space and proceeds in the next step to grasp it.
In a high-level robot, an optical means (such as a television system) is
used to scan the background of the object. It recognizes the pattern and
determines the presence and orientation of the object. A computer is necessary to process signals in the pattern-recognition process (see Figure 11-2).
Feedback signal
Input device
Output device
Figure 11-2. Robot Using a Pattern Recording Process.
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In some applications, the computerized robot recognizes the presence and
orientation of each mechanical part by a pattern recognition process that
consists of reading the code numbers attached to it. Then the robot picks up the
part and moves it to an appropriate place for assembling, and there it
assembles several parts into a component. A well-programmed digital computer
acts as a controller.
Temperature Control System. Figure 11-3 shows a schematic diagram of
temperature control of an electric furnace. The temperature in the electric
furnace is measured by a thermometer, which is an analog device. The analog
temperature is converted to a digital temperature by an A/D converter. The
digital temperature is fed to a controller through an interface. This digital
temperature is compared with the programmed input temperature, and if
there is any discrepancy (error), the controller sends out a signal to the
heater, through an interface, amplifier, and relay, to bring the furnace
temperature to a desired value.
Temperature Control of the Passenger Compartment of a Car. Figure
11-4 shows a functional diagram of temperature control of the passenger
compartment of a car. The desired temperature, converted to a voltage, is the
input to the controller. The actual temperature of the passenger compartment
is converted to voltage through a sensor and is fed back to the controller for
comparison with the input. The ambient temperature and radiation heat
transfer from the sun, which are not constant while the car is driven, act as
disturbances. This system employs both feedback control and feed forward
C ontrol . ( Feed forward control gives corrective action before the
affect the output.)
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Figure 11-4. Temperature Control of Passenger Compartment of a Car.
The controller receives the input signal, output signal, and signals from
sensors from disturbance sources. The controller sends out an optimal control
signal to the air conditioner to control the amount of cooling air so that the
passenger compartment temperature is equal to the desired temperature.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1.The second paragraph mainly discusses .......…. .
a. the factors affecting the speed of an engine
b. the disturbances created due to fluctuations in the speed of an
c. the mechanism of the speed governor to adjust th speed of an
d. the rate of oil flow in a speed governor to reduce external disturbances
2. In the speed control system just described, the amount of fuel to be
applied to the engine is known as .......…. .
a. the force
b. the disturbance
c. the feedback signal
d. the actuating signal
3. According to the text, robots ........…. .
a. cannot work at very low temperatures
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b. cannot handle jobs without error
c. are equipped with proper devices to be able to perform the tasks
d. are powered to handle various jobs performed by man in industry
4. A high-level robot ........…. .
a. performs pattern recognition to determine how to pick up, move,
and assemble the parts into a component
b. determines the presence and orientation of parts to change their
shapes and weights if needed
c. does not have enough control over the process it goes through
d. does not have the proper program for pattern recognition in its
computer memory
5. The examples given in this text are of ........…… control systems.
a. closed-loop
b. open-loop
c. hydraulic
d. pneumatic
B. Write the answers to the following questions.
1. What is the function of a speed governor?
2. How freely is an industrial robot able to move?
3. What is the use of a microswitch installed in the arm of a robot?
4. How does a temperature control system work?
5. What are considered disturbances in the temperature control of the
passenger compartment of a car?
6. What constitutes the temperature control of the passenger compartment of a car?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Historical Review
The first significant work in automatic control was James Watt’s centrifugal
governor for the speed control of a steam engine in the eighteenth century.
Other significant works in the early stages of development of control theory
were due to Minorsky, Hazen, and Nyquist, among many others. In 1922,
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Minorsky worked on automatic controllers for steering ships and showed how
stability could be determined from the differential equations describing the
system. In 1932, Nyquist developed a relatively simple procedure for
determining the stability of closed-loop systems on the basis of open-loop
response to steady-state sinusoidal inputs. In 1934 Hazen, who introduced the
term servomechanisms for position control systems, discussed the design of
relay servomechanisms capable of closely following a changing input.
During the decade of the 1940s, frequency-response methods made it
possible for engineers to design linear closed-loop control systems that
satisfied performance requirements. From the end of the 1940s to early 1950s,
the root-locus method due to Evans was fully developed.
The frequency-response and root-locus methods, which are the core of
classical control theory, lead to systems that are stable and satisfy a set of
more or less arbitrary performance requirements. Such systems are, in general,
acceptable but not optimal in any meaningful sense. Since the late 1950s, the
emphasis in control design problems has been shifted from the design of one
of many systems that work to the design of one optimal system in some
meaningful sense.
As modern plants with many inputs and outputs become more and more
complex, the description of a modern control system requires a large number
of equations. Classical control theory, which deals only with single-input,
single-output systems, becomes powerless for multiple-input, multiple-output
systems. Since about 1960, because the availability of digital computers made
possible time-domain analysis of complex systems, modern control theory,
based on time-domain analysis and synthesis using state variables, has been
developed to cope with the increased complexity of modern plants and the
stringent requirements on accuracy, weight, and cost in military, space, and
industrial applications.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1. actuate
2. adaptive control system
3. aircraft-autopiloting system
4. alternate biasing
5. automatic regulating system
6. classical control theory
7. closed-loop control system
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9. complex system
10. electromagnetic valve
11. external disturbance
12. feedback control
13. internal disturbance
14. linear control system
15. microswitch
16. missile-guidance system
17. multiple-input
18. multiple-output
19. numerical control
20. open-loop system
21. pneumatic control
22. position control system
23. preset
24. process control system
25. robot locus method
26. robotic system
27. sensory device
28. servomechanism
29. servo system
30. space-vehicle system
31. steady-state
32. time-domain analysis
33. Watt’s centrifugal governor
34. Watt’s (Watt’s flyball governor) governor
35. Watt’s speed governor
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Section One: Reading Comprehension
Magnetic Line Voltage Starters
Magnetic control means the use of electromagnetic energy to close switches.
Line voltage (across the line) magnetic starters arc electromechanical devices
that provide a sale, convenient, and economic means of full voltage starting
and stopping motors. In addition, these devices can be controlled remotely.
They are used when a full-voltage starling torque may be applied safely to the
driven machinery and when the current inrush resulting from across-the-line
starting is not objectionable to the power system. Control for these starters is
usually provided by pilot devices such as push buttons, float switches, timing
relays, etc. Automatic control is obtained from the use of some of these pilot
Magnetic vs Manual Starters
Using manual control, the starter must be mounted so that it is easily within
reach of the machine operator. With magnetic control, push-button stations
are mounted nearby, but automatic control pilot devices can be mounted
almost anywhere on the machine. The push buttons and automatic pilot
devices can be connected by control wiring into the coil circuit of a remotely
mounted starter, possibly closer to the motor to shorten the power circuit.
In the construction of a magnetic controller, the armature is mechanically connected to a set of contacts so that, when the armature moves to its
closed position, the contracts also close. There are different variations and
positions, but the operating principle is the same.
The simple up-and-down motion of a solenoid-operated, three-pole
magnetic switch is shown in Figure 12-1. Not shown are the motor overload
relays and maintaining and auxiliary electrical contacts. Double break contacts
are used on this type of starter to cut the voltage in half on each contact, thus
providing high arc rupturing capacity and longer contact life.
The operating principle that makes a magnetic starter different from a
manual starter is the use of an electromagnet. Electrical control equipment
makes extensive use of a device called a solenoid. This electromechanical
device is used to operate motor starters, contactors, relays, and valves. By
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Figure 12-1. Three-Pole, Solenoid-Operated Magnetic Switch (Contactor) and
Electrical Wiring Symbols.
placing a coil of many turns of wire around a soft iron core, the magnetic flux
set up by the energized coil tends to be concentrated; therefore, the magnetic
field effect is strengthened. Since the iron core is the path of least resistance
to the magnetic lines of force, magnetic attraction concentrates according to
the shape of the magnet core.
There are several different variations in design of the basic solenoid
magnetic core and coil. Figure 12-2 shows a few examples. As shown in the
solenoid design of Figure 12-2C, linkage to the movable contacts assembly is
obtained through a hole in the movable plunger. The plunger is shown in the
open de-energized position.
The center leg of each of the E-shaped magnet cores in Figures 12-2B
and C is ground shorter than the outside legs to prevent the magnetic switch
from accidentally staying closed (due to residual magnetism) when power is
Figure 12-3 shows a manufactured magnet structure and how the starter
contacts are mounted on the armature.
When a magnetic motor starter coil is energized and the armature has
been sealed in, it is held tightly against the magnet assembly. A small air gap
is always deliberately placed in the center leg, iron circuit. When the coil is
de-energized, a small amount of magnetism remains.If it were not for this gap
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Figure 12-2. Some Variations of Basic Magnet Core and Coil Configurations of
in the iron circuit, the residual magnetism might be enough to hold the
movable armature in the sealed-in position. This knowledge can be important
to the electrician when troubleshooting a motor that will not stop.
12-3. Magnet Structure (Left) and Movable Contacts and Armature
Guide Assembly (Right) of a Four-Pole Magnetic Switch (Courtesy Square D Co.).
The OFF or OPEN position is obtained by de-energizing the coil and
allowing the force of gravity or spring tension to release the plunger from the
magnet body, thereby opening the electrical contacts. The actual contact
surfaces of the plunger and core body are machine ground to insure a high
degree of flatness on the contact surfaces so that operation on alternating
current is quieter. Improper alignment of the contacting surfaces and foreign
matter between the surfaces may cause a noisy hum on alternating-current
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Another source of noise is loose laminations. The magnet body and
plunger (armature) are made up of thin sheets of iron laminated and riveted
together to reduce eddy currents and hysteresis, iron losses showing up as heal
Figure 12-4. Types of Magnet Cores.
(see Figure 12-4). Eddy currents are shorted currents induced in the metal by
the transformer action of an ac coil. Although these currents are small, they
heal up the metal, create an iron loss, and contribute to inefficiency. At one
time, laminations in magnets were insulated from each other by a thin,
nonmagnetic coating; however, it was found that the normal oxidation of the
metallic laminations reduces the effects of eddy currents to a satisfactory
degree, thus eliminating the need for a coating.
Part I. Comprehension Exercises
A. Put "T" for true and "F" for false statements. Justify your
........ 1. Line voltage magnetic starters provide electric motors with safe
........ 2. Magnetic controllers can be mounted anywhere on the machine.
........ 3. The armature plays the major part in the construction of a
magnetic controller.
........ 4. A solenoid is an electromechanical device.
........ 5. Alternating-current magnets are inclined to create noise.
........ 6. Eddy currents do not affect the efficiency of alternating-current
B. Choose a, b, c, or d which best completes each item.
1. The first paragraph mainly discusses .......... .
a. the application of full voltage starting torque to the driven mach
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b. the application of safe current to power systems
c. the use of electromagnetic energy to control electrical machinery
d. the use of pilot devices to control magnetic starters
2. In solenoid, magnetic concentration …….. .
a. directly depends on the design of the solenoid magnetic core and coil
b. directly depends on the number of turns of wire around the solenoid
c. reduces the amount of core resistance to the magnetic lines o f force
d. reduces the magnetic field effect caused by the magnetic coil
3. The shorter length of the center leg of an E-shaped magnet core …..…..
when power is disconnected.
a. concentrates more energy than the outside legs
b. concentrates less energy than the outside legs
c. causes the magnetic switch to stay closed
d. causes the magnetic switch to function properly
4. As we understand from the text, ………. .
a. the shape of the magnet core of a solenoid determines the
concentration of the magnetic attraction
b. the magnetic flux set up by the coil concentrates according to the
energy applied to it
c. manual starters are practically more useful than magnetic starters
d. solenoids are not widely used in electrical control equipment
5. If a small air gap were not placed in the center leg of the magnet core,
………. when the magnetic motor starter coil was de-energized.
a. the residual magnetism would increase
b. the plunger might release the magnet body
c. a motor might not stop working
d. a motor might not continue working
C. Answer the following questions orally.
1. How are magnetic starters usually controlled?
2. What is the function of double break contacts on a solenoid- operand,
three-pole magnetic switch?
3. What is the operating principle of a magnetic controller?
4. What is the use of the hole in the movable plunger?
5. How do the electrical contacts of a magnetic starter open?
6. What are eddy currents?
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Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. An electric controller , …… , is used to start and stop a motor.
a. a switch
b. a starter
c. a push button
d. a timing relay
2. An ………. starter connects the motor to the supply without the use of a
resistance or autotransformer to reduce the voltage.
a. across-the-line
b. automatic
c. autotransformer
d. increment
3. Due to……….,ferromagnetic bodies retain a certain magnetization after
the magnetizing force has been removed.
a. solenoid conductivity
b. solenoid resistivity
c. residual modulation
d. residual magnetism
4. Voltage induced in the body of a conducting mass by a variation of
magnetic flux results in ………. .
a. eddy currents
b. electromagnetic energy
c. hysteresis
d. magnetism
5. Alternating-current magnets are laminated and riveted together to
reduce ...……. .
a. magnetic flux
c. induction
b. magnetic charge
d. hysteresis
B. Fill in the blanks with the appropriate form of the
1. Energize
a. The available ……... is the amount of work that a system is capable of
b. When excessive current is drawn, the relay de- ………. the starter and
stops the motor.
2. Devise
a. The human operator may be replaced by a mechanical ,electrical ,or
similar ………. .
b. Circuit controllers are ………. to close and open electric.
3. Maintain
a. The controller functions to ……… the furnace temperature close to
the varying set point.
b. Quick tripping may be …….. by a time delay overload relay.
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4. Vary
a. Most physical systems are nonlinear to ………. extents.
b. If the range of……of the system variables is not wide, the system may
be linearized within a relatively small range of variation of variables.
c. In a continuous time control system, all system …….. are functions
of a continuous time t.
d. A time-invariant control system is one whose parameters do not
……….with time.
5. Overload
a. To provide .. ……. or running protection to keep a motor from
overheating, overload relays are used on starters to limit the amount
of current drawn to a predetermined value.
b. Unless the cause of the ………. has been removed, theoverload relay
will trip again.
c. A magnetic overload relay is used to stop an electrical machine when
it is ……. .
C. Fill in the blanks with the following words.
cut off
Instantaneous trip current relays are ……….to take a motor off the ………
as soon as a predetermined load condition is ……... For example, when a
blockage of ……..on a woodworking machine causes a sudden..……. current,
an instantaneous trip relay can the motor quickly. After the cause of the
……… is removed, the motor can be restarted ………because the relay resets
itself as soon as the ……… is removed.
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. As a result, the coil of the magnetic relay must be wound with wire
large enough in size to pass the motor current.
b. In some cases, the relay may also be used so that it is actuated when the
current falls to a certain value.
c. The magnetic overload relay coil is connected in series directly with the
motor or is indirectly connected by current transformers (as in circuits
with large motors).
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d. They are used when an electrical contact must be opened or closed as
the actuating current rises to a certain value.
e. These overload relays operate by current intensity and not heat.
Section Two: Further Reading
Motor Overheat
An electric motor does not know enough to quit when the load gets too much
for it. It keeps going until it burns out. If a motor is subjected, over a period
of time, to internal or external heat levels that are high enough to destroy the
insulation on the motor windings, it will fail-burn out.
A solution to this problem might be to install a larger motor whose
capacity is in excess of the normal horsepower required. This is not too
practical since there are other reasons for a motor to overheat besides excess
loads. A motor will run cooler in the winter cold than in the summer heat of
a tropical climate. A high, surrounding air temperature (ambient temperature)
has the same effect as higher-than-normal current flow through a motor-it
tends to deteriorate the insulation on the motor windings.
High ambient temperature is also created by poor ventilation of the
motor. Motors must get rid of their heat, so any obstructions to this process
must be avoided. High inrush currents of excessive starting create heat within
the motor. The same is true with starting heavy loads. There are several other
related causes that generate heat within a motor such as voltage unbalance,
low voltage, and single phasing. In addition, when the rotating member of the
motor will not turn (a condition called locked rotor), heat is generated.
The ideal overload protection for a motor is an element with current
sensing properties very similar to the heating curve of the motor. This would
act to open the motor circuit when full load is exceeded. The operation of the
protective device is ideal if the motor is allowed to carry small, short and
harmless overloads, but is quickly disconnected from the line when an
overload has persisted too long. Dual element, or time-delay, fuses may
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provide motor overload protection, but they have the disadvantage of being
nonrenewable and must be replaced.
An overload relay is added to the magnetic switch that was shown in
Figure 12-1. Now it is called a motor starter. The overload relay assembly is
the heart of motor protection. The motor can do no more work than the
overload relay permits. Like the dual element fuse, the overload relay has
characteristics permitting it to hold in during the motor accelerating period
when the inrush current is drawn. Nevertheless, it still provides protection on
small overloads above full-load current when the motor is running. Unlike the
fuse, the overload relay can be reset. It can withstand repeated trip and reset
cycles without need of replacement. It is emphasized that the overload relay
does not provide short circuit protection. This is the function of overcurrent
protective equipment like fuses and circuit breakers, generally located in the
disconnecting switch enclosure.
Current drawn by a motor is a convenient and accurate measure of the
motor load and motor heating. Therefore, the device used for overload
protection, the overload relay, is usually connected with the motor current. It
is provided as part of the starter or controller. As the relay carries the motor
current, it is affected by that current. If a dangerous over-current condition
occurs, it operates or trips the relay to open the control circuit of the
magnetic starter and disconnect the motor from the line; this helps insure the
maximum operating life of the motor. In a manual starter, an overload trips a
mechanical latch causing the starter contacts to open and disconnect the motor
from the line.
The controller is normally installed in the same room or area as the
motor. This makes it subject to the same ambient temperature as the motor.
The tripping characteristic of the proper thermal overload relay will then be
affected by room temperature exactly as the motor is affected. This is done by
selecting a thermal relay element (from a chart provided by the manufacturer)
that trips at the danger temperature for .the motor windings. When excessive
current is drawn, the relay de-energizes the starter and stops the motor.
Overload relays can be classified as being either thermal or magnetic.
Magnetic overload relays react only to current excesses, and are not affected
by temperature. As the name implies, thermal overload relays depend on the
rising ambient temperature and temperatures caused by the overload current to
trip the overload mechanism.
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Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. Overloading causes a motor to overheat which in turn results in ……. .
a. starter breakdown
b. its failure
c. low coil currents
d. low magnetic pull
2. As we understand from the text, a motor can best be protected against
overheating by………. .
a. an overload relay
b. reducing the ambient temperature
c. a time-delay fuse
d. reducing high inrush currents
3. Paragraph Jive mainly describes………. .
a. a motor accelerating period
b. a motor running period
c. the function of a dual element fuse
d. the function of an overload relay
4. It may be inferred from the text that ……… .
a. it is impossible to design a motor that will adjust itself to all the
various changes of heat
b. it is possible to design a motor that will adjust itself to all the various
changes of heat
c. poor ventilation does not always create high ambient temperature
d. starting heavy loads do not always create heat within a motor
5. It is true that ……… .
a. the overload relay has nothing to do with the motor current
b. the controller is usually installed in a different room away from the
c. current drawn by a motor is directly proportional to the motor load
d. current drawn by a motor de-energizes the starter and stops the motor
B. Write the answers to the following questions.
1. Why are time-delay fuses not considered the best overload protective
2. How dose an overload protective device function?
3. What is a dual element fuse?
4. What is a motor starter?
5. What factors may cause a motor to be overheated?
6. Why is an overload relay connected with the motor current?
7. How are overload relays classified?
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Section Three: Translation Activities
A. Translate the following passage into Persian.
Time Limit Overload Relays
Time delay overload relays make use of the oil dash pot principle. Motor
current passing through the coil of the relay exerts a magnetic pull on a
plunger. The magnetic flux set up inside the coil tends to raise the plunger
which is attached to a piston immersed in oil. As the current increases in the
relay coil, so does the magnetic flux. The force of gravity is overcome and the
plunger and piston move upward. During this upward movement, oil is forced
through bypass holes in the piston. As a result, the operation of the contacts
is delayed. A valve disc is turned to open or close bypass holes of various sizes
in the piston. This action changes the rate of oil flow and so adjusts the time
delay factor. The rate of upward travel-of the core and piston-depends
directly upon the degree of overload. The greater the current load, the faster
the upward movement. As the rate of upward movement increases, the relay
tripping time decreases.
This inverse time characteristic prevents the relay from tripping on the
normal starting current or on harmless momentary overloads. In these cases,
the line current drops to its normal value before the operating coil is able to
lift the core and piston far enough to operate the overload control contacts.
However, if the overcurrent continues for a prolonged period, the core is
pulled far enough to operate the contacts. As the line current increases, the
relay tripping time decreases. Tripping current adjustment is achieved by
adjusting the plunger core with respect to the overload relay coil. Quick
tripping is obtained through the use of a light trade dashpot oil and by
adjustment of the oil bypass holes.
A valve in the piston allows almost instantaneous resetting of the circuit
to restart the motor. The current must then be reduced to a very low value
before the relay will reset. This action is accomplished automatically when the
tripping of the relay disconnects the motor from the line. Magnetic overload
relays are available with either automatic reset contacts or hand reset
B. Find the Persian equivalents of the following terms
expressions and write them in the spaces provided.
1. armature
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2. autotransformer
3. auxiliary
4. blockage
5. contactor
6. contribute
7. eddy current
8. electromechanical float switch
9. excessive current
10. increment
11. intensity
12. laminated core
13. lock-rotor
14. machinery
15. magnetic control
16. magnetic flux
17. manual
18. movable plunger
19. nonlinear
20. poor ventilation
21. residual magnetism
22. residual modulation
23. solid core
24. time eddy fuse
25. time-invariant control system
26. timing relay
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Figure 13-3. Schematic Symbols for the Phototransistor. the Photodiode, and the
When the photo transistor is in darkness, no electrons are emitted by the base
junction, and the transistor is turned off. When the photo transistor is in the
presence of light, it turns on and permits current to flow through the relay
coil. The diode connected parallel to the relay coil is known as a kickback or
freewheeling diode. Its function is to prevent induced voltage spikes from
occurring when the current suddenly stops flowing through the coil and the
magnetic field collapses.
In the circuit shown in Figure 13-4, the relay coil will turn on when the
photo transistor is in the presence of light, and turn off when the phototransistor is in darkness. Some circuits may require the reverse operation. This
can be accomplished by adding a resistor and a junction transistor to the
circuit, Figure 13-5. In this circuit a common junction transistor is used to
Figure 13-4. Photo transistor Controls Relay Coil.
control the current flow through the relay coil. Resistor R 1 limits the current
flow through the base of the junction transistor. When the photo transistor is
in darkness, it has a very high resistance. This permits current to flow to the
base of the junction transistor and turn it on. When the photo transistor is in
the presence of light, it turns on and connects the base of the junction
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transistor to the negative side of the battery, This causes the junction
transistor to turn off. The photo transistor in the circuit is used as a stealer
transistor. A stealer transistor steals the base current away from some other
transistor to keep it turned off.
Figure 13-5. The Relay Turns on When the Photo transistor Is in Darkness.
Some circuits may require the photo transistor to have a higher gain than
it has under normal conditions. This can be accomplished by using the
photo transistor as the driver for a Darlington amplifier circuit, Figure 13-6. A
Darlington amplifier circuit generally has a gain of over 10,000.
Photo diodes and photo-SCRS are used in circuits similar to those shown
for the photo transistor. The photo diode will permit current to flow through it
in the presence of light. The photo-SCR has the same operating characteristics as a common junction SCR. The only difference is that light is used to
trigger the gate when using a photo-SCR.
Figure 13-6. The Phototransistor Is Used as the Driver for a Darlington Amplifier.
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Regardless of the type of photoemissive device used, or the type circuit it is
used in, the greatest advantage of the photoemissive device is speed. A
photoemissive device can turn on or off in a few microseconds. Photovoltaic or
photoconductive devices generally require several milliseconds to turn on or off.
This makes the use of photoemissive devices imperative in high speed switching
Photoconductive Devices
Photoconductive devices exhibit a change of resistance due to the presence or
absence of light. The most common photoconductive device is the cadmium sulfide
cell or cad cell. The cad cell has a resistance of about 50 ohms in direct sunlight
and several hundred thousand ohms in darkness. It is generally used as a light
sensitive switch. The schematic symbol for a cad cell is shown in Figure 13-7, Figure
13-8 shows a typical cad cell.
Figure 13-9. Schematic Symbol for a
cad cell
Figure 13-8. Cad Cell (Courtesy EG &
G Vactec , Inc.)
Figure 13-9 shows a basic circuit of a cad cell being used to control a
relay. When the cad cell is in darkness, its resistance is high. This prevents the
amount of current needed to turn the relay on from flowing through the
circuit. When the cad cell is in the presence of light, its resistance is low. The
amount of current needed to operate the relay can now flow through the
Although this circuit will work if the cad cell is large enough to handle
the current, it has a couple of problems.
1.There is no way to adjust the sensitivity of the circuit. Photo-operated
switches are generally located in many different areas of a plant. The
surrounding light intensity can vary from one area to another. It is,
therefore, necessary to be able to adjust the sensor for the amount of
light needed to operate it.
2.The sense of operation of the circuit cannot be changed. The circuit
shown in Figure 13-9 permits the relay to turn on when the cad cell is
in the presence of light. There may be conditions that would make it
desirable to turn the relay on when the cad cell is in darkness.
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Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements . Justify your
……. 1. Photodetectors sense the presence of an object by physical contact
with the object.
……. 2. Photovoltaic devices release electrons in the presence of light.
……. 3. The amount of voltage produced by a solar cell depends on the
material it is made of.
……. 4. A large and a small solar cell made of the same material produce
the same amount of voltage.
……. 5. Photoemissive devices are of three types.
……. 6. A phototransistor may not emit electrons in darkness.
……. 7. Photovoltaic devices turn on and off faster than photoemissive
B. Choose a, b, c, or d which best completes each item.
1. As we understand from the text, the increasing application of photodetectors in industry is due to their .......... .
a. speed and their ability to sense the presence or absence of almost
any object
b. flexibility of being used in almost every type of industry
c. ability to sense objects without making any contact with the objects
d. all of the above
2. Paragraph three mainly discusses .......... solar cells.
a. the amount of voltage produced by
b. the amount of current produced by
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c. the characteristics of
d. the characteristics of the material used in
3. In a circuit consisting of a resistor, a junction transistor, a diode, a
photo transistor, and a relay …….. .
a. the relay coil will turn on when the photo transistor is in the
presence of light
b. the relay coil will turn on when the photo transistor is in darkness
c. the photo transistor exhibits a high resistance in the presence of light
d. the photo transistor emits more electrons in the presence of light
4. The junction transistor used in the circuit controls the current flow
through………. .
a. the relay coil
b. the resistor
c. the diode
d. the phototransistor
5. It is true that a cad cell ……… .
a. is not usually used as a light sensitive switch
b. is usually used to produce high currents
c. acts much faster than photovoltaic cell
d. exhibits lower resistance in the presence of sunlight
C. Answer the following questions orally.
1. What does the amount of current produced by a solar cell depend on?
2. How can 1he amount of voltage produced by solar cells be increased?
3. What are the three categories of photo-operated devices?
4. What is the difference between photo voltaic and photoemissive
5. What is the function of the kickback diode in the circuit?
6. What are the problems of a cad cell circuit?
7. What is the advantage of photovoltaic cells over other photo-operated
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. Photodiode current is very low. If this current is to be used effectively
in control applications, it must be amplified by an external current
amplifier or by a device called
a. a photodetector
b. a phototransistor
c. a photodiode
d. a photoconductor
2. The collector-base junction of the phototransistor acts as
a. a transistor
b. a resistor
c. a photodetector
d. a photodiode
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3. Solar cells often called ……… devices are made of silicon and have the
ability to produce a voltage in the presence of light.
a. photovoltaic
b. photoemissive
c. photodiode
d. photoconductive
4. Current developed by a phototransistor is dependent mainly on ………
of light and very little on the applied voltage.
a. the frequency
b. the wavelength
c. the intensity
d. both a and b
5. The ability of the low-current gate circuit of ……… to control large
amounts of power in its anode-cathode circuit makes this device
particularly useful in industrial electronics.
a. photo-SCR
b. photodetector
c. photodiode
d. phototransistor
B. Fill in the blanks with the appropriate form of the words
1. Industry
a. Photodetectors designed for ……… use are made to be mounted and
used in different ways.
b. Optoelectronics deals with light-sensitive semiconductor devices used
in …….. .
2. Dope
a. Photovoltaic cells consist of a single semiconductor crystal which has
been ……… with both N- and P-type materials.
b. Photodiodes and phototransistors are each sensitive to a specific range
of light frequencies. What these frequencies are, is determined by the
materials from which they are constructed and by the extent
of ………….. of the junctions.
c. When light falls on the PN junction, which is the boundary of these
………., a voltage appears across the junction.
3. Exhibit
a. Photodiodes and phototransistors…….. unique spectral characteristics
and have a variety of uses in industry.
b. Selective properties within the range of the visible spectrum are
……….by the average human eye.
c. Photoconductive cells ..…….. the particular property that their
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resistance decreases in the presence of light and increases in the
absence of light.
4. Contact
a. The switch .......... of relays may be normally open or normally closed.
b. The relay coil and .......... terminals can usually be located by
5. Sense
a. The current that a photovoltaic cell can deliver to a load depends on
the area of the light- .......... material which makes up the cell.
b. The .......….. of a photocell to a particular color depends on the
nature of the material from which the cell is constructed, and on the
manner of its construction,
c. Manufacturers, specify the spectral …......... of their devices in the form
of a frequency response curve.
C. Fill in the blanks with the following words.
Light behaves like an electromagnetic radiation, and ….......... with light are
the characteristics of wavelength and ........….. . The wavelength of light
determines the ........…. of that light. White light consists of many ........... which
may be separated by a prism. The .......... eye responds to the prismatic colors
ranging from .......... to red . The wavelengths of light which the .......... can ‘see’
are in the nanometer .......... of 400 nm to 700 nm. A .......... is a 10 9 part of a
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. A glass window in the housing permits light to fall on the active
material of the cell.
b. Photoconductive cells are made of a thin layer of semiconductor
material such as cadmium sulfide, cadmium selenide, or lead sulfide.
c. The ceil simply acts as a conductor whose resistance changes when
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d. The semiconductor layer is enclosed in a sealed housing.
e. Photoconductive cells exhibit the peculiar property that their resistance
decreases in the presence of light and increases in the absence of light.
Section Two: Further Reading
Optoelectronics-The Phototransistor
There is a variety of semiconductor-junction light-sensitive devices which fall
under the heading of optoelectronic devices. These devices are being used, in
conjunction with other semiconductors, in such diverse control applications as
automatic light level controls in photocopy machines to computer control of
machine tools.
Included in the optoelectronics’ family are: light-emitting diodes (LEDs),
photodiodes, phototransistors, photo-SCRs [also called light-activated SCRs
(LASCRs)], optocouplers or optoisolators, and solid-state relays (SSRs).
Silicon photodiodes are light-sensitive devices, also called photodetectors,
which convert light signals into electrical signals. A window or lens permits
light to fall on the junction (Figure 13-10). When light shines on the
reversebiased PN photodiode junction, hole-electron pairs are created. The
Figure 13-10. Light Falls on the PN Junction of a Photodiode.
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movement of these hole-electron pairs in a properly connected circuit results
in current flow. The current is proportional to the intensity of light and is also
affected by the frequency of the light falling on the photojunction,
The response of the human eye is not uniform in the range of the visible
spectrum. As Figure 13-11 shows, the eye is most sensitive to light whose
wavelength is 550 nm and falls off to 400 and 700 nm. The spectral response
of the eye, then, is 400 to 700 nm, peaking at 550 nm.
Figure 13-11. Spectral Response of the Human Eye. (To convert wavelength to angstrom
units, multiply above values by 10.)
The spectral response of a silicon photodiode is shown in Figure 13-12.
We see that maximum sensitivity is to radiation at 900 nm, and that the total
response is in the range from 400 to 1100 nm. This includes response both in
Figure 13-12. Spectral Response of a Silicon Photodiode (Motorola).
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the visual range and outside it. We call the total range ‘light’ even though
light normally refers to the frequencies within the visual(human)spectrum. The
spectral response curve of a particular silicon photodiode depends on the
geometry and the extent of doping of the junction. It is apparent, then, for
maximum efficiency, that the spectral characteristics of the light source (light
emitter) used with a photodiode must match the characteristics of the
photodiode. So in the case of the photodiode in Figure 13-12, the light source for
it must have a wavelength close to 900 nm.
The current developed by a photodiode is very low. This current cannot be
used directly in control applications but must be amplified. After amplification, the photocurrent may be high enough to be used in a control system
for example, to set a relay. The phototransistor is a light detector which
combines a photodiode and a transistor amplifier. Figure 13-13 shows an NPN
phototransistor. Here a lens concentrates the light on a very thin P-type
wafer, sandwiched in between an N-type collector and an emitter. Although
the phototransistor has three sections, only two leads may issue from the
housing, the emitter and collector leads. In this device base current is supplied
Figure 13-13. NPN Phototransistor
by the current created by the light falling on the base-collector photodiode
junction. Some phototransistors have a third lead issuing from the housing. In
such a phototransistor, base bias is provided from an external circuit, on which
the photodiode current is superimposed.
Current in a phototransistor is dependent mainly on the intensity of light
entering the transistor window and is little affected by the voltage applied to
the external circuit. Figure 13-14 is a graph of collector current I c , as a
function of collector-emitter voltage VCE and as a function of illumination H.
It is evident that the phototransistor acts as a constant-current source, and
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Figure 13-14. Collector Characteristic for the MRD300 (Motorola).
Figure 13-15.
Spectral Responses of a Phototransistor and the Human Eye
(General Electric).
Figure 13-15 is the frequency response curve of a phototransistor for
comparison with the response of the human eye. Peak sensitivity of the
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phototransistor is at about 900 nm.
The angular alignment of a phototransistor and its source of illumination are important considerations. The reason is that the illumination of the
photojunction is proportional to the cosine of the angle between the direction
of radiation of the light source and the perpendicular to the surface of the
photojunction. In addition, the optical lens, or window, and its size further
affect the response of the phototransistor.
Phototransistor Relay
Figure 13-16 is the circuit diagram of a one-stage relay employing an NPN
120 V 60 HZ
Figure 13-16. One-Stage Phototransistor-Operated Relay.
Figure 13-17. Tow-Stage Phototransistor and Amplifier.
phototransistor Q,. A very sensitive relay is connected in the collector circuit
and is actuated when a strong light is focused directly on the transistor lens.
The resistance and the pickup current of the relay must be relatively low so
that the transistor may be operated within its rated characteristics. A is a
milliammeter used for measuring collector current. If an ordinary 60-W light
bulb is used as the light source, the light must be held fairly close to the
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photowindow to actuate the relay.
Less light is required to actuate the same relay in the circuit shown in
Figure 13-17, Hence the light source need not be held so close to the light
window, In this circuit, the output of phototransistor Ql is amplified by
transistor Q2" The relay coil is connected in the collector of Q2' Photocurrent
flow in Ql through R1, biases Q2' In the absence of photocurrent (that is, in
the absence of light on QJ, Q2 is biased to cutoff, since the emitter and the
base of Q2 are at the same potential. When light shines on Ql emitter current
flows through Rl and forward-biases the emitter to base of Q2' thus causing
current to flow in Q2' If the light is strong enough, there will be sufficient
current in the collector of Q2 to actuate the relay. As in the preceding circuit,
a sensitive relay is required.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. The third paragraph mainly describes …….. .
a. how the intensity of light affects current
b. how the frequency of light affects current
c. the movement of hole-electron pairs
d. the mechanism of photodiodes
2. As we understand from the text, ……… .
a. the spectral characteristics of various photodiodes are different
b. the spectral characteristics of all photodiodes are similar
c. the visual spectrum extends outside the spectral response of a
d. the efficiently of a photodiode does not relate to the light source
3. The higher the doping of a particular photodiode junction ………. .
a. the higher the number of hole-electron pairs created
b. the lower the number of hole-electron pairs created
c. the higher the efficiency of its spectral response
d. the lower the efficiency of its spectral response
4. It is true that …….. .
a. the photodiode current is amplified to be used in a control system
b. the primary source of the photodiode current is an external circuit
c. a photo transistor consists of two transistor amplifiers
d. a photo transistor consists of a collector and an emitter
5. The last two paragraphs mainly discuss
a. relay circuits with one or two transistors
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b. milliammeters used in different relay circuits
c. the kinds of relays used in a circuit
d. the number of transistors employed in a relay circuit
6. We may infer from the text that
a. an optoelectronic device is an electronic device containing optic and
electric ports
b. an optoelectronic device cannot be used as a control element in
complicated machines
c. the current developed by a photodiode is high enough to be used
directly in control applications
d. the spectral response of the eye can be extended to the spectral
response of a silicon photodiode
B. Write the answers to the following questions.
1. What are some applications of optoelectronic devices?
2. What is a photodetector?
3. How does a photodiode create current?
4. What is the current developed by a photodiode proportional to?
5. How does the spectral response of the eye differ from that of a
6. What is known as light?
7. What brings about the maximum efficiency of the spectral response of a
particular photodiode?
8. What does a phototransistor consist of?
9. What is the illumination of the photojunction proportional to?
10. How do you describe the mechanism of a relay circuit having two
Section Three: Translation Activities
A. Translate the following passage into Persian.
Photovoltaic Cells
These are also light-sensitive semiconductor devices, but they produce a
voltage when illuminated, which increases as the intensity of light falling on
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the semiconductor junction of this two-element cell increases. The usual basic
material from which these cells are made today is silicon or selenium.
Photovoltaic cells convert light into electric energy, which may be used
directly to supply small amounts of electric power for electrically powered
devices. Because of the low levels of power which photovoltaic cells generate,
they have been used in the past in low-power devices such as light meters and
photographic exposure meters. However, with an improvement in the
efficiency of these cells, more power has been produced, as in solar cells,
which are photovoltaic devices. Solar cells appear destined to play a substantial role in the development of new sources of energy.
Photovoltaic cells consist of a single semiconductor crystal which has
been doped with both N- and P-type materials. When light tails on the PN
junction, which is the boundary of these dopants, a voltage appears across the
junction. About 0.6 V is developed by the photovoltaic cell in bright sunlight.
The amount of power the cell can deliver depends on the extent of its active
surface. An average cell will produce about 30 milliwatts per square inch (30
m W/in 2 ) of surface, operating into a load of 4 Ω. To increase the power
output, large banks of cells are used in series and parallel combinations. An
example is the use of solar cells to power the experimental circuits on lunar
and space modules.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1 . angular alignment
2 . collapse
3 . Darlington amplifier circuit
4 . dope
5 . evident
6 . exhibit
7 . lexibility
8 . illumination
9 . light activated SCR (LASCR)
10 . light emitting diode (LED)
11 . optocoupler
12 . optoelectronic
13 . optoisolator
14 . perpendicular
15 . photocell
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16. photodetector
17. photodiode
18. photoemissive
19. photojunction
20. photo-operated device
21. photo-operated switch
22. photovoltaic
23. photowindow
24. reverse biased
25. semiconductor
26. solid-state relay (SSR)
27. spectrum
28. stealer transistor
29. trigger
30. wafer
31. wavelength
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Unit 14
ng Comprehension
rument Classification and
e subdivided into separate classes according to several
classifications are useful in broadly establishing several
cular instruments such as accuracy, cost, and general
rent applications.
vided into active or passive ones according to whether the
is entirely produced by the quantity being measured or
ity being measured simply modulates the magnitude of
source. This is illustrated by examples.
assive instrument is the pressure-measuring device shown
e pressure of the fluid is translated into a movement of a
ale. The energy expended in moving the pointer is derived
change in pressure measured: there are no other energy
assive Pressure Gauge.
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Figure 14-2. Petrol-Tank Level Indicator.
active instrument is a float-type petrol-tank level indicator
gure 14-2. Here, the change in petrol level moves a
, and the output signal consists of a proportion of the
urce applied across the two ends of the potentiometer. The
put signal comes from the external power source: the
float system is merely modulating the value of the voltage
ower source.
ts, the external power source is usually in electrical form,
it can be other forms of energy such as a pneumatic or
t difference between active and passive instruments is the
ent resolution which can be obtained. With the simple
own, the amount of movement made by the pointer for a
change is closely defined by the nature of the instrument.
to increase measurement resolution by making the pointer
he pointer tip moves through a longer arc, the scope for
is clearly limited by the practical limit of how long the
ently be. In an active instrument, however, adjustment of the
external energy input allows much greater control over
Type Instruments
just mentioned is a good example of a deflection type of
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instrument, where the value of the quantity being measured is displayed in
terms of the amount of movement of a pointer.
An alternative type of pressure gauge is the dead-weight gauge shown in
Figure 14-3 which is a null-type instrument. Here, weights are put on top of
the piston until the downward force balances the fluid pressure. Weights are
added until the piston reaches a dat um level, known as the null point.
Pressure measurement is made in terms of the value of the weights needed to
reach this null position.
Figure 14-3. Dead-Weight Pressure Gauge.
The accuracy of these two instruments depends on different things. For
the first one, it depends on the linearity and calibration of the spring, while
for the second, it relies on the calibration of the weights. As calibration of
weights is much easier than careful choice and calibration of a linearcharacteristic spring, this means that the second type of instrument will
normally be the more accurate. This is in accordance with the general rule
that null-type instruments are more accurate than deflection types.
Monitoring/Control Instruments
An important distinction between different instruments is whether they are
suitable only for monitoring functions or whether their output is in a form that
can be directly included as part of an automatic control system. Instruments
which only give an audio or visual indication of the magnitude of the physical
quantity measured, such as a liquid-in-glass thermometer, are only suitable
for monitoring purposes. This class normally includes all null-type
instruments and most passive transducers.
For an instrument to be suitable for inclusion in an automatic control
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system, its output must be in a suitable form for direct input into the
controller. Usually, this means that an instrument with an electrical output is
required, although other forms of output such as optical or pneumatic signals
are used in some systems.
Analog/Digital Instruments
An analog instrument gives an output which varies continuously as the
quantity being measured changes. The output can have an infinite number of
values within the range that the instrument is designed to measure. The
deflection type of pressure gauge is a good example of an analog instrument.
As the input value changes, the pointer moves with a smooth continuous
motion. While the pointer can therefore be in an infinite number of positions
within its range of movement, the number of different positions which the eye
can discriminate between is strictly limited, this discrimination being dependent upon how large the scale is and how finely it is divided.
A digital instrument has an output which varies in discrete steps and so
can only have a finite number of values. The rev-counter sketched in Figure
14-4 is an example of a digital instrument. A cam is attached to the revolving
body whose motion is being measured, and on each revolution the cam opens
and closes a switch. The switching operations are counted by an electronic
counter. This system can only count whole revolutions and cannot discriminate any motion which is less than a full revolution.
Figure 14-4. Rev-Counter.
The distinction between analog and digital instruments has become
particularly important with the rapid growth in the application of micro-
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computers to automatic control systems. Any digital computer system, of
which the microcomputer is but one example, performs its computations in
digital form. An instrument whose output is in digital form is therefore
particularly advantageous in such applications, as it can be interfaced directly
to the control computer. Analog instruments must be interfaced to the
microcomputer by an analog-to-digital (ND) converter, which converts the
analog output signal from the instrument into an equivalent digital quantity
which can be read into the computer.
Part I. Comprehension Exercises
A. Put "T" for true and "F" for false statements. Justify your
……… 1. The output of a passive instrument is directly produced by the
quantity being measured.
……… 2. The output of an active instrument is determined by an external
power source.
……… 3. In passive instruments, the external power source can be various
forms of energy.
……… 4. In a passive pressure gauge, the energy used to move the pointer is
derived from the change in pressure of the fluid measured.
……… 5. A deflection-type instrument is more accurate than a null-type
……… 6. An analog computer does its calculations one step at a time
whereas a digital computer continuously works out its calculation
B. Choose a, b, c, or d which best completes each item.
1. The function of the float system in a petrol-tank level indicator is …… .
a. to modulate the value of the voltage from the external power source
b. to evaluate the amount of energy from the external power source
c. to calculate the amount of the external voltage applied across the
two ends of the potentiometer
d. to adjust the energy produced by the quantity being measured
2. It is understood from the text that the level of measurement resolution
obtained by …….. .
a. passive instruments can be highly under control
b. active instruments can be highly under control
c. a simple pressure gauge cannot be increased
d. a simple pressure gauge can be highly increased
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3. It may be inferred from the text that ……… .
a. the scope for improving measurement resolution is infinite
b. the scope for improving me4surement resolution is not infinite
c. active and passive instruments can both be used for accurate
d. active and passive instruments do not vary in their basic structures
4. The mechanism of a null-type instrument is based on …….. .
a. the force applied by weights
b. the linearity of a spring
c. the movement of a pointer
d. the state of equilibrium
5. It is true that active instruments ……….. .
a. are only suitable for monitoring purposes
b. are not as effective for measuring purposes as passive ones
c. can be used for control purposes
d. cannot be used for control purposes
6. It can be inferred from the last paragraph that ………. .
a. the time involved in the process of converting an analog signal to a
digital quantity can be critical in the control of fast processes
b. the time involved in the process of converting an analog signal to a
digital quantity is too small to be considered a disadvantage of the
analog instrument
c. the analog to digital converter does not affect the speed of operation
of the control computer
d. the analog to digital converter does not affect the accuracy of the
control computer
C. Answer the following questions orally.
1. How do you describe the mechanism of a passive pressure gauge?
2. What constitutes the potentiometer in the petrol-tank level indicator?
3. How are null-type instruments different from deflection types?
4. Why is an analog instrument interfaced to a microcomputer by an AID
5. How does a rev-counter work?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. A passive transducer…….
a. can adjust the magnitude of some external power source
b. can transform the energy obtained from an external source
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c. has no source of power other than the input signals
d. may have either internal or external sources of power
2. A three-terminal rheostat, or a resistor with one or more adjustable
sliding contacts, that functions as an adjustable voltage divider is called
………. .
a. a potentiometer
b. a converter
c. a gauge
d. a counter
3. The energy in the output signal of ………... is entirely produced by the
quantity being measured.
a. a passive pressure gauge
b. a petrol-tank level indicator
c. a deflection type instrument
d. a digital instrument
4. A ………. instrument can be directly included as part of an automatic
control system.
a. null-type
b. passive
c. control
d. monitoring
5. The energy in the output signal of ………. instruments comes from an
external source of power.
a. passive
b. active
c. monitoring
d. null-type
B. Fill in the blanks with the appropriate form of the words
1. Require
a. Choice between active and passive instruments for a particular
application involves carefully balancing the measurement-resolution
………. against cost.
b. The higher the gate current, the lower the anode-cathode voltage
………. to turn the SCR on.
c. The current passing through the coil …….. a definite time interval to
reach its maximum or steady-state value.
d. A zener diode may be used as a voltage regulator for a load …….. a
voltage equal to the zener voltage.
2. Obtain
a. Angular velocity measurements can be ……… by differentiating the
output signal from angular displacement transducers.
b. In measurement systems which contain an angular acceleration
transducer, such as a gyro accelerometer, it is possible to…….. a
velocity measurement by integrating the acceleration measurement
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3. Consist
a. The drag-cup tachometer has a central spindle carrying a permanent
magnet which rotates inside a non-magnetic drag-cup ……….. of a
cylindrical sleeve of electrically conductive material.
b. A bistable multivibrator ………. of two direct cross-coupled dc
4. Accurate
a. The drag-cup tachometer has a typical measurement ………of  0.5%
and is commonly used in the speedometers of motor vehicles and as a
speed indicator for aero engines.
b. If the input data entering the computer are correct and if the
program of instructions is reliable, then we can expect that the
computer generally will produce ………… output.
5. Use
a. A silicon rectifier, when ………. to convert ac to dc, acts as a closed
switch when its anode is positive relative to its cathode and as an
open switch when its anode is negative relative to its cathode.
b. The vertical amplifiers of an oscilloscope must be calibrated if the
scope is to be ………. for measuring the amplitude of waveforms.
c. The ability of the low-current gate circuit of an SCR to control large
amounts of power in its anode-cathode circuit makes this device
particularly ………. in industrial electronics.
C. Fill in the blanks with the following words.
Probably the most common form of ………. output device used is the dc
……... .This is a relatively simple device which …….. speeds up to about 5000
rpm ……….. an accuracy of  l %. Where better ………. is required within a
similar range of……… measurement, ac tachometers are used. The …… rotor
type has an accuracy of  0.25% and ……… rotor types have accuracies up to
 0.05%.
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. It consists of a pair of spherical balls pivoted on a rotating shaft.
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b. The pointer can be arranged to give a visual indication of speed by
causing it to move in front of a calibrated scale, or its motion can be
converted by a translational displacement transducer into an electrical
c. These balls move outward under the influence of centrifugal forces as
the rotational velocity of the shaft increases and lift a pointer against
the resistance of a spring.
d. The mechanical flyball is a velocity-measuring instrument which was
first developed many years ago and is still used extensively in speedgoverning systems for engines, turbines, etc.
Section Two: Further Reading
Rotational Velocity Measurement
The main application of rotational velocity transducers is in speed control
systems. They also provide the usual means of measuring translational
velocities, which are transformed into rotational motions for measurement
purposes by suitable gearing. Many different instruments and techniques are
available for measuring rotational velocity as presented below.
DC Tachometric Generator
The dc tachometric generator, or dc tachometer as it is generally known, has
an output which is approximately proportional to its speed of rotation. Its
basic structure is identical to that found in a standard dc generator used for
producing power, and is shown in Figure 14-5. Both permanent-magnet types
and separately excited field types are used. However, certain aspects of the
designare optimized to improve its accuracy as a speed-measuring instrument.
One significant design modification is to reduce the weight of the rotor by
constructing the windings on a hollow fiberglass shell. The effect of this is to
minimize any loading effect of the instrument on the system being measured.
The dc output voltage from the instrument is of a relatively high
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Permanent magnet
Wound rotor supplied
via slip rings
Figure 14-5. The DC Tachometer.
magnitude, giving a high measurement sensitivity which is typically 5 V per
1000 rpm. The direction of rotation is determined by the polarity of the
output voltage. A common range of measurement is 0-5000 rpm Maximum
non-linearity is usually about  1 % of the full-scale reading.
One problem with these devices which can cause problems under some
circumstances is the presence of an ac ripple in the output signal. The
magnitude of this can be up to 2% of the output dc level.
AC Tachometric Generator
The ac tachometric generator, or ac tachometer as it is generally known, has
an output approximately proportional to rotational speed, as in the dc
tachogenerator. Its mechanical structure takes the form of a two-phase
induction motor, with two stator windings and (usually) a drag-cup rotor, as
shown in Figure 14-6. One of the stator windings is excited with an ac voltage
and the measurement signal is taken from the output voltage induced in the
second winding. The magnitude of this output voltage is zero when the rotor
is stationary, and otherwise proportional to the angular velocity of the rotor.
The direction of rotation is determined by the phase of the output voltage,
which switches by 1800 as the direction reverses. Therefore, both the phase
and magnitude of the output voltage have to be measured. A typical range of
measurement is 0-4000 rpm with an accuracy of  0.05% of full-scale reading
While the form of ac tachometer described above is the commonest one,
a second form also exists. This has a squirrel-cage rather than a drag-cup rotor
and so is cheaper. Its structure and mode of operation are otherwise identical,
but the measurement accuracy is reduced.
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Figure 14-6. The A C Tachometer.
Variable-Reluctance Velocity Transducer
The form of a variable-reluctance transducer is shown in Figure 14-7. It can be
seen that this consists of two parts, a rotating disk connected to the moving
Figure 14-7. Variable-Reluctance Transducer.
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body being measured and a pick-up unit. The rotating disk is constructed from
a bonded-fiber material into which soft iron poles are inserted at regular
intervals around its periphery. The pick-up unit consists of a permanent
magnet with a shaped pole piece which carries a wound coil. The distance
between the pick-up and the outer perimeter of the disk is around 0.5 mm.
As the disk rotates, the soft iron inserts on the disk move in turn past the
pick-up unit. As each iron insert moves toward the pole piece, the
reluctance of the magnetic circuit increases and hence the flux in the pole
piece also increases. Similarly, the flux in the pole piece decreases as each
iron insert moves away from the pick-up unit, and the pattern of flux changes
with time as shown in Figure 14-8. The changing magnetic flux inside the pickup coil causes a voltage to be induced in the coil whose magnitude is proportional to the rate of change of flux. This voltage is positive while the flux is
increasing and negative while it is decreasing and therefore its variation with
lime is as shown in Figure 14-9.
Figure 14-8. Pattern of Flux Change With Rotation of a Variable-Reluctance Transducer.
Figure 14-9. Pattern of Induced Voltage Change With Rotation of a VariableReluctance Transducer.
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The form of this output can be regarded as a sequence of positive and
negative pulses whose frequency is proportional to the rotational velocity of the
disk. This can be converted into an analog, varying-amplitude, dc voltage output
by means of a frequency-to-voltage converter circuit connected to the output
terminals of the pick-up. However, greater measurement accuracy can be
obtained by converting the output waveform into sharp pulses which are
counted by an electronic counter. It is normal procedure to produce the pulses
at each instant that the induced voltage in the coil changes sign as it passes
through zero. This is achieved by electronic means.
The maximum angular velocity which the instrument can measure is
limited because of the finite width of the induced pulses. As the velocity
increases, the distance between the pulses is reduced, and at a certain velocity the
pulses start to overlap. At this point, the pulse counter ceases to be able to
distinguish the separate pulses.
The total pulse count measured over a certain length of time only gives
information about the average velocity over that period. Measurement of the
actual velocities at the instants of time that each output pulse occurs can be
achieved by the scheme shown in Figure 14-10. In this circuit, the pulses from the
transducer gate the train of pulses from a I MHz clock into a counter, Control
logic resets the counter and updates the digital output value after receipt of
each pulse from the transducer. The measurement resolution of this system is
highest when the speed of rotation is low.
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Photoelectric Pulse-Counting Methods
An alternative to the variable-reluctance transducer, b u t which uses very similar
principles to it, is the method where pulses are produced by photoelectric
techniques and counted. These pulses are generated by one of the two alternate
methods illustrated in Figure 14-11. In Figure 14-ll(a), the pulses are produced as
the windows in a slotted disk pass in sequence between a light source and a
detector. The alternate form, Figure 14-H(b), has both light source and
detector mounted on the same side of a reflective disk which has black sectors
painted onto it at regular angular intervals. In cither case, the pulses are
counted by an electronic counter. The frequency of the pulses is proportional to
the angular velocity of the body connected to the rotating disk. Pulses
generated in this manner are narrower than those generated by magnetic
means and so the instrument is capable of measuring higher velocities.
Figure 14-10. Scheme to Measure Instantaneous Angular Velocities.
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5. How can an analog, varying amplitude, dc voltage be obtained from a
variable-reluctance transducer?
6. What does the finite width of the induced pulses in a variablereluctance transducer lead to?
7. How can actual velocities at the instants of time that each output pulse
occurs be measured?
8. How do you describe the photoelectric pulse generation techniques?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Stroboscopic Methods
The stroboscopic technique of rotational velocity measurement operates on a
similar physical principle to the variable-reluctance and photoelectric pulsecounting methods, except that the pluses involved consist of flashes of light
generated electronically and whose frequency is adjustable so that it can be
matched with the frequency of occurrence of some feature on the rotating
body being measured. This feature can either be some naturally occurring one
such as the spokes of a wheel or gear teeth, or it can be an artificially created
pattern of black and white stripes. In either case, the rotating body appears
stationary when the frequencies of the light pulses and body features are in
synchronism. Flashing rates up to 25,000 per minute are available from
commercial stroboscopes, according to the range of velocity measurement
required, and the typical measurement accuracy obtained is  1% of the
Measurement of the flashing rate at which the rotating body appears
stationary does not automatically indicate the rotational velocity, because
synchronism also occurs when the flashing rate is some integral submultiple of
the rotational speed. The practical procedure followed is therefore to adjust
the flashing rate until synchronism is obtained at the largest flashing rate
possible, R 1 . The flashing rate is then carefully decreased until synchronism is
again achieved at the next lower flashing rate, R 2 The rotational velocity is
then given by:
R1 R2
R1  R2
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B. Find the Persian equivalents of the following terms
expressions and write them in the spaces provided.
1. beam splitter
2. calibration
3. cross-coupled
4. deflection type
5. discriminate
6. drag-cup
8. feature
9. fiber optic
10. gyroscope
11. interferometer
12. monitoring control instrument
13. null type
14. photoelectric pulse-counting method
15. pilot
16. potentiometer
17. rotational velocity
18. sleeve
19. spherical
20. spindle
21. squirrel-cage
22. stationary core
23. stroboscopic method
24. tachogenerator
25. tachometer
26. variable-reluctance transducer
27. zener diode
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Section One: Reading Comprehension
Television Broadcasting
The term ‘broadcasting’ means to send out in all directions. As illustrated in
Figure 16-1, the transmitting antenna radiates electromagnetic radio waves
that can be picked up by the receiving antenna. For commercial television
broadcast stations, the service area is about 25 to 75 mi in all directions from
the transmitter. The radiation is in the form of two rf carrier waves,
modulated by the desired information. Amplitude modulation(AM) is used
for the picture signal. However, frequency modulation (FM) is used for the
sound signal.
Referring to Figure 16-1, we see that the desired sound for the televised
program is converted by the microphone to an audio signal, which is amplified
for the sound-signal transmitter. For transmission of the picture, the camera
tube converts the visual information into electrical signal variations. A camera
tube is a cathode-ray tube with a photoelectric image plate.
The electrical variations from the camera tube become the video signal,
which contains the desired picture information. The video signal is amplified
and coupled to the picture-signal transmitter for broadcasting to receivers in
the service area.
Separate carrier waves are used for the picture signal and sound signal,
but they are radiated by one transmitting antenna. Furthermore, the picture
and sound signals are included in the broadcast channel for each station. A
television channel for a commercial broadcast station is made 6 MHz wide to
include both the picture and sound. At the receiver also, one antenna is used
for the picture and sound signals.
The receiving antenna intercepts the radiated picture and sound carrier
signals, which are then amplified and detected in the receiver. The detector
output includes the desired video signal containing the information needed to
reproduce the picture. Then the recovered video signal is amplified and
coupled to a picture lube that converts the electric signal back into light.
Reproducing the Picture
The picture tube is very similar to the cathode-ray tube used in the
oscilloscope. The glass envelope contains an electron-gun structure that
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Figure 16-1. Block Diagram of the Television Broadcasting System.
produces a beam of electrons aimed at the fluorescent screen. When the
electron beam strikes the screen, it emits light.
When the signal voltage makes the control grid less negative, the beam
current is increased, making the spot of light on the screen brighter. More
negative grid voltage reduces the brightness. If the grid voltage is negative
enough to cut off the electron-beam current at the picture tube, there will be
no light. This state corresponds to black. A color picture tube has three
electron guns for the tricolor screen.
The picture tube is also called a kinescope or a CRT. Its function is to
convert the video signal into a picture.
Scanning and Synchronizing
In order for the camera tube to convert the picture information into video
signal, the image is dissected into a series of horizontal lines. Similarly, the
picture tube reassembles the image line by line. These horizontal lines are
produced by making the electron beam scan across the screen. There are 525
lines per picture frame. In addition to this horizontal scanning, vertical
scanning is necessary to spread the lines from top to bottom of the screen.
There are 30 complete picture frames per second.
Furthermore, the scanning at the camera tube and picture tube must be
synchronized, or timed, with respect to the video signal. The synchronization
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is necessary to reassemble the picture information on the correct lines. These
functions are provided by the block of scanning and synchronizing circuits
shown in Figure 16-1 for the transmitter and receiver. The term ‘synchronizing’ is usually abbreviated sync.
Most programs are produced live in the studio but recorded on video
tape at a convenient time to be shown later. The quality is so good that the
picture looks practically the same as a live program. The studio also has
projectors to use 35-mm still slides, opaque slides, and motion-picture film,
either 16 or 35 mm, as the program source.
For remote pickups, as in broadcasting a sports event, the signal is
relayed to the studio for broadcasting in the assigned channel. When there is a
national hookup for important programs, each station in the network receives
video signal by intercity relay links, usually leased from the telephone
company. A system for satellite relay stations covering the country is being
developed for this nationwide television service.
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
……1. Radio waves are used to carry the picture and the sound signals.
……2. Carrier waves are modulated by the desired information after being
picked up by the receiving antenna.
……3. Sound signals are converted to audio signals prior to transmission.
……4. The radiated picture and sound carrier signals are amplified and
detected in the receiver.
……5. The electron beam produced by the electron gun strikes the
fluorescent screen causing the screen to emit light.
……6. The electron beam scans the screen horizontally to produce picture
B. Choose a, b, c, or d which best completes each item.
1. As we understand from the first paragraph, …….. .
a. amplitude modulation may also be used for audio frequency signals
b. audio waves are modulated and then broadcast to the receiver
c. carrier waves are used to produce sound and picture signals
d. sound and picture signals cannot be transmitted directly
2. According to the text, …….. .
a. visual information is converted into electrical signal variations
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b. electrical variations become video signals containing the desired
picture information
c. video signals are amplified and broadcast to the receiver
d. all of the above
3. Picture and sound signals are …….. .
a. radiated by different transmitting antennas
b. received by two receiving antennas.
c. transmitted through the same television channel
d. broadcast by the same carrier wave
4. At the receiver, ……… .
a. the detector produces the desired video signal
b. the video signal produced by the detector is amplified and coupled to
a picture tube
c. the picture tube converts the video signal into a picture
d. all of the above
5. As we understand from the text, ……… .
a. the picture tube used in television sets has a different mechanism
from that used in an oscilloscope
b. the sound carrier signal is detected in a receiver different from that
detecting the picture carrier signal
c. by varying the negative potential on the grid in the electron gun, the
intensity of the beam varies
d. by increasing the negative potential on the grid the brightness of the
spot of light on the screen increases
B. Answer the following questions orally
1. How are picture and sound signals transmitted to the receiving antenna?
2. What is the function of the camera tube?
3. What is the function of the picture tube?
4. How many picture frames are produced on the screen per second?
5. How does synchronization affect a picture frame?
6. How are remote pickups done?
7. What are intercity relay links?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. In amplitude modulation, ……….. of a carrier signal is varied by the
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modulating voltage, whose frequency is invariably lower than that of
the carrier.
a. the frequency
b. the amplitude
c. the phase
d. all of the above
2. A cathode-ray …….. is an electron-beam tube in which the beam can be
focused on to a small cross section on a luminescent screen and varied
in position and intensity to produce a visible pattern.
a. instrument
b. oscillograph
c. tube
d. oscilloscope
3. An ……….. is an electrode structure that produces one or more electron
a. electron gun
b. electron tube
c. electronic controller
d. electronic converter
4. A cathode-ray tube used to produce an image by variation of the beam
intensity as the beam scans a raster is called ………. .
a. an electron gun
b. an oscillator
c. a picture element
d. a picture tube
5. A ……….. is an electron tube used to provide an image in color by the
scanning of a raster and by varying the intensity of excitation of
phosphors to produce light of the chosen primary colors.
a. color-purity magnet
b. color-picture tube
c. color-selecting-electrode system d. color-field collector
B. Fill in the blanks with the appropriate form of
1. Amplify
a. Diodes may be combined with electronic DC …….. to form an
electronic voltmeter or other electronic instruments.
b. The horizontal amplifier in an oscilloscope ………. the time-base
voltage from the sweep oscillator, providing a control for the width
of the resulting pattern.
c. The operational amplifier serves as the heart of the analog computer,
because it possesses the widely useful ability to provide a high value
of precisely controlled …….. .
2. Transmit
a. In an AM ………… , amplitude modulation can be generated at any
point after the radio frequency source.
b. There are several different systems for TV ………and reception.
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c. The video chain at the ……… station begins with a transducer which
converts light into electric signals.
3. Radiate
a. Any power escaping into free space is governed by the characteristics
of free space. If such power escapes on purpose, it is said to have
been ……… .
b. Free space is the space that does not interfere with the normal ……….
and propagation of radio waves.
c. Antennas ………. electromagnetic waves, or, putting it differently,
radiation will result from the flow of high-frequency current in a
suitable circuit.
4. Detect
a. The diode is used for AM demodulation or ……... .
b. A ……… detects the presence of electric waves.
c. A radar ………. the presence of objects and their distance by the
transmission and return of electromagnetic energy.
5. Synchronize
a. The task of the ……… circuits in a television receiver is to process
received information, in such a way as to ensure that the vertical and
horizontal oscillators in the receiver work at the correct frequencies.
b. In television, the synchronizing signal is employed for the ……… of
c. In a …….. computer, each event, or the performance of each
operation, starts as a result of a signal generated by a clock.
C. Fill in the blanks with the following words.
When practical radio transmission started in the year 1901, the ……….
radio frequencies of about 100 kHz were used for………distances of hundreds
to thousands of miles. As radio ………, higher frequencies were used for
services requiring more………. . Now we have television broadcasting in the
VHF ………. of 30 to 300 MHz and the UHF band of 300 to 3,000 MHz.
However, the…….. for wireless transmission becomes much shorter at these
high ………… .Broadcasting is practically limited to the line-of-sight distance
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.......... the transmitting and receiving antennas in the VHF and UHF bands.
The .......... service range is up to 75 mi for VHF ..........and 25 to 35 mi for
UHF stations.
D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. Among the more important factors are the modulation system
used, the operating frequency, the operating range, and the type of
display required which in turn depends on the destination of the
intelligence received .
b. Receivers range widely in complexity, from a very simple crystal
receiver with headphones, to a far more complex radar receiver with its
involved antenna arrangements and visual display system.
c. A great variety of receivers are used in communication systems, because
the exact form of a suitable receiver for a particular usage is influenced
by a number of different factors .
d. Both these processes are the reverse of the corresponding transmitter
processes .
e. Whatever the receiver, its most important function is demodulation
( and sometimes also decoding).
Section Two: Further Reading
The Cathode Ray Tube
The cathode ray tube operates as follows. First, electrons are emitted from a
heated cathode. Then these electrons are accelerated to give them a high
velocity. Next they are formed into a beam which can be deflected vertically and
horizontally. Finally they are made to strike a screen coated on its inner surface
with a phosphor.
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The CRT comprises an electron gun and a deflection system enclosed in
a glass tube with a phosphor coated screen. The electron gun forms the
electrons into a beam. It contains a cathode which is heated to produce a
stream of electrons. On the same axis as the cathode is a cylinder known as
the grid. By varying the negative potential on the grid, the intensity of the
beam can be varied. A system of three anodes follows. These accelerate the
beam and also operate as a lens to focus the beam on the screen as a small
dot. Varying the potential on the central anode, a 2 , allows the focus to be
On leaving the electron gun, the beam passes through the deflection
system. There are two systems in use today for moving the electron beam
around on the face of the CRT. They are the electrostatic and the
electromagnetic deflection systems. The type of deflection will depend on the
ultimate use of the CRT display. Electrostatic deflection is accomplished by
placing four metal plates inside the neck of the CRT with connecting wires to
the outside for voltage application, Two of the plates control the vertical
movement of the beam, and two plates control the horizontal movement.
The two vertical plates are placed above and below the path of the
electron beam, as shown in Figure 16-2, while the two horizontal plates are on
either side of the electron beam path. When one plate is grounded and a
positive AC voltage is applied to the other plate, the path of the electron
beam will bend toward the positive plate. The beam is moving too fast to
strike the deflection plate, so only its direction is changed. The advantage of
electrostatic deflection is that it performs equally well at all frequencies, from
DC to the limits of the CRT's ability to produce a trace. The disadvantage is
the small deflection angles of about ±15°. A 21 in. picture tube would be
about 48 in. long.
Path of the electron beam
Figure 16-2. Electrostatic Deflection Plates.
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Electromagnetic deflection uses two sets of coils in place of the deflection plates. A sawtooth current is required through the coils in order to
achieve linear trace and retrace. The coils, called the deflection yoke, fit
around the outside of the neck of the CRT and have an inductive component
as well as the wire resistance considered to be in series with the inductance. A
sawtooth voltage applied to the resistive element will produce a sawtooth
current through the resistance. However, a square wave voltage must be
applied to the inductive portion of the yoke to produce a sawtooth current
through the inductance. The sum of the square wave voltage and the sawtooth
voltage is called a trapezoidal voltage waveshape. The advantage of the
electromagnetic deflection system is the wide deflection angles of about ±60° The
disadvantage is that coils work best at only one frequency. However, this
is not really a disadvantage because in television, each half of the yoke is
required to work at only one frequency-60 Hz for the vertical and 15,750 Hz
for the horizontal.
The vertical circuit consists of an oscillator that will operate near 60
cycles to generate the voltage wave needed for deflection. The vertical
frequency adjustment (called the vertical hold control) is part of the oscillator
circuit. This is a free-running oscillator that will be synchronized by the
station signal at the exact frequency of that station. The vertical size control
(or height control) adjusts the DC voltage to the oscillator and controls the
strength of the output signal, which determines the size of the vertical scan of
the CRT (see Figure 16-3).
Figure 16-3. Vertical Deflection Block Diagram.
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The oscillator signal is amplified by the vertical output power amplifier
to convert the voltage signal to a current signal to drive the deflection yoke.
The vertical linearity control is generally a part of this circuit. The linearity
control is a gain control that varies the operating Q point of the amplifier and
distorts the waveshape by compressing the top or bottom of the wave as
required to achieve the best distribution of the vertical scan from top to
bottom of the CRT screen. The size control and the linearity controls interact
for best picture results.
Vertical yoke currents of 1 A are common. The vertical oscillator and
the vertical amplifier with the deflection yoke comprise the entire vertical
The horizontal deflection section has the same functions as the vertical
system, but is developed in a slightly different manner. The horizontal
oscillator used for television was the first practical use of the phase-locked
loop. The oscillator has a free-running frequency of 15,750 cycles adjustable
by a control setting (horizontal hold control), as shown in Figure 16-4. That is,
with no input signal, the circuit components are selected to control the
frequency. The oscillator is a voltage-controlled circuit that depends on the
DC voltage at the input to establish the frequency of operation.
Figure 16-4. Horizontal Deflection and High-Voltage Block Diagram.
A feedback signal from the horizontal circuit is applied to a phase
detector, which compares the frequency of an incoming pulse from the station
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(when a station is selected) and develops the DC voltage to control the
oscillator. The oscillator output is amplified by a power amplifier to develop
the current signal used to drive the horizontal section of the deflection yoke.
This exchange from voltage 10 current is accomplished through a horizontal
output transformer. Because of the higher horizontal frequency (compared to
the vertical), the horizontal yoke current need not be as great as the current
in the vertical deflection system.
The final element is the phosphor coated screen. When the electron
beam strikes the screen, the phosphor coating fluoresces. Various colors of
light are produced depending on the phosphor used.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. The intensity of the beam is controlled by varying …….. .
a. the number of electrons emitted for the cathode
b. the potential on the central anode
c. the negative potential on the grid
d. the form of the electron beam
2. The second paragraph mainly discusses the mechanism of ......… .
a. the electron gun
b. the cathode ray tube
c. the deflection system
d. the central anode
3. The electrostatic deflection system consists of ………… .
a. four plates and a screen
b. two sets of deflection plates
c. an electron gun and a screen
d. four plates and an electron gun
4. It is true that ………. .
a. electrostatic deflection system uses two sets of coils
b. electromagnetic deflection system uses two sets of plates
c. electromagnetic deflection system works well at all frequencies
d. electrostatic deflection system performs well at all frequencies
5. According to the text, ……… .
a. the electromagnetic deflection system is made up of the
deflection yoke which contains inductive and resistive elements
b. the sawtooth voltage applied to the resistive element produces a
sawtooth current through the inductance
c. either a sawtooth voltage or a square wave voltage is required to
produce a trapezoidal voltage waveshape
d. both a sawtooth current and a square wave voltage are needed to
produce a trapezoidal waveshape
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6.The voltage wave required for the vertical deflection is provided by
……….. .
a. an oscillator
b. an inductor
c. a resistor
d. a converter
7. The best picture results from the interaction between …….. .
a. the vertical hold control and the oscillator
b. the oscillator and the vertical size control
c. the size control and the linearity controls
d. the amplifier and the oscillator
B. Write the answers to the following questions.
1. What is the source of electrons for the electron beam?
2. What is the function of the electron gun?
3. In what way is the system of anodes like a lens?
4. What are the two deflection systems called?
5. How do the plates change the direction of the beam?
6. What is a trapezoidal voltage waveshape?
7. What does the vertical circuit consist of?
8. How does the phase detector in the horizontal deflection system
control the oscillator?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Color Television
The block diagram in Figure 16-1illustrates the television broadcasting system
for monochrome. In color television, a color camera is necessary at the
transmitter and a color picture tube at the receiver. The color
provides video signal for the red, green, and blue picture
information. A color picture tube has red, green, and blue phosphors on the
reproduce the picture in color. A Typical color picture tube has three electron
guns for the tricolor screen. The phosphors can be dot trios of red, green, and
blue, or vertical stripes of color. Then each gun produces an electron beam to
illuminate the red, green, or blue phosphor dots on the fluorescent screen.
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Although the camera and picture tube operate with red, green, and blue
all other colors including white can be reproduced by combinations of these
three colors. Furthermore, in commercial television, the red, green, and blue
signals are combined for broadcasting. The purpose is to transmit only a
chrominance signal for color and a luminance signal that contains the
monochrome information. It is necessary to transmit the luminance signal so
that monochrome receivers can reproduce the picture in black and white. The
chrominance signal or chroma signal has all the information needed to
reproduce the picture in color.
The luminance signal is called the Y video signal. The chrominance
signal can be called the C signal. Actually, the C signal is a modulated
subcarrier of 3.58 MHz. This 3.58-MHz C signal modulates the assigned
picture carrier in the standard 6-MHz television broadcast channel.
Furthermore, the 3.58-MHz chrominance signal itself is modulated by two
color video signals. The process of interweaving the Y signal for luminance
with the 3.58-MHz color subcarrier signal for color is called multiplexing. In
terms of the modulated chrominance signal, 3.58 MHz is the frequency for
color in the television broadcast system.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
2. deflection yoke
3. detect
4. distort
5. free-running oscillator
6. hookup
7. illuminate
8. intercept
9. kinescope
10. modulate
12. sawtooth current
13. subcarrier
14. trapezoidal
15. vertical hold control
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Section One: Reading Comprehension
Transmission Lines
Transmission lines are a means of conveying signals or power from one point
to another. From such a broad definition, any system of wires can be
considered as forming one or more transmission lines. However, if the
properties of these lines must be taken into account, the lines might as well
be arranged in some simple, constant pattern. This will make the properties
much easier to calculate, and it will also make them constant for any type of
transmission line. Thus all practical transmission lines are arranged in some
uniform pattern: this simplifies calculations, reduces costs, and increases
convenience. There are two types of commonly used transmission lines. The
parallel-wire (balanced) line is shown in Figure 17-lb, and the coaxial
(unbalanced) line in Figure 17-la.
Inner conductor
(a) Coaxial (unbalanced) line
Outer casing
(b) Parallel wire (balanced) line
Figure 17-1. Transmission Lines.
The parallel-wire line is employed where balanced properties are required: for instance, in connecting a folded-dipole antenna to a TV receiver or
a rhombic antenna to an HF transmitter. On the other hand, the coaxial line
is used when unbalanced properties are needed, as in the interconnection of a
broadcast transmitter to its grounded antenna. It is also employed at UHF
and microwave frequencies, to avoid the risk of radiation from the
transmission line itself.
Any system of conductors is likely to radiate if the conductor separation
approaches a half-wavelength at the operating frequency. This is far more
likely to occur in a parallel-wire line than in a coaxial line, whose outer
conductor surrounds the inner one and is invariably grounded. Accordingly,
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parallel-wire lines are never used for microwaves, whereas coaxial lines may be
employed for frequencies up to at least 18 GHz. However, from the general
point of view the limit is on the lowest usable frequency; below about 1 GHz,
waveguide cross-sectional dimensions become inconveniently large. Within
each broad grouping or type of transmission line there is an astonishing
variety of different kinds, dictated by various applications. Lines may be rigid
or flexible, air-spaced or filled with different dielectrics, with smooth or
corrugated conductors as the circumstances warrant. Different diameters and
properties are also available. Flexible lines are naturally more convenient than
rigid ones, since they may be bent to follow any physical layout and are much
easier to stow and transport. On the other hand, rigid cables can generally
carry much higher powers, and it is easier to make them air-dielectric rather
than filled with a solid dielectric. This consideration is important, especially
for high powers, since all solid dielectrics have significantly higher losses than
air, particularly as frequencies are increased.
Rigid coaxial air-dielectric lines consist of an inner and outer conductor
with spacers of low-loss dielectric separating the two every few centimeters.
There may be a sheath around the outer conductor to prevent corrosion, but
this is not always the case. A flexible air-dielectric cable generally has
corrugations in both the inner and the outer conductor, running at right
angles to its length, and a spiral of dielectric material between the two.
The power-handling ability of a transmission line is limited by flashover
between the conductors due to a high-voltage gradient breaking down the
dielectric. It depends on the type of dielectric material used, as well as the
distance between the conductors. Thus, for the high-power cables employed in
transmitters, nitrogen under pressure may be used to fill the cable and reduce
flashover. Since nitrogen is less reactive than the oxygen component of air,
corrosion is reduced as well. Dry air under pressure is also used as a means of
keeping out moisture. Clearly, as the power transmitted is increased, so must
be the cross-sectional dimensions of the cable.
Since each conductor has a certain length and diameter, it must have
resistance and inductance; since there are two wires close to each other, there
must be capacitance between them. Finally, the wires are separated by a
medium called the dielectric, which cannot be perfect in its insulation; the
current leakage through it can be represented by a shunt conductance. The
resulting equivalent circuit is as shown in Figure 17-2.
At radio frequencies, the inductive reactance is much larger than the
resistance. The capacitive susceptance is also much larger than the shunt
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Figure 17-2. General Equivalent Circuit of Transmission Line.
conductance. Thus both R and G may be ignored, resulting in a line that is
considered lossless (as a very good approximation for RF calculations). The
equivalent circuit is simplified as shown in Figure 17-3.
It is to be noted that the quantities L, R, C, and G, shown in Figures
17-2 and 17-3, are all measured per unit length, e.g., per meter, because they
Occur continuously along the line.
They are thus distributed throughout
the length of the line. Under no
circumstances can they be assumed
to be lumped at any one point.
Figure 17-3. Transmission-Line RF Equivalent
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
........ 1. The parallel-wire line may be used to connect a broadcast
transmitter to its grounded antenna.
........ 2. The parallel-wire line is usually used at HF and UHF frequencies.
........ 3. A parallel-wire line is more liable to radiation than a coaxial line.
........ 4. The higher the frequency, the higher the power loss a solid
dielectric will have.
........ 5. The sheath around the outer conductor of a rigid coaxial
air-dielectric cable is not of much use.
B. Choose a, b, c, or d which best completes each item.
1.The first paragraph mainly discusses ........….. .
a. the basic principles of transmission lines
b. the basic calculations for transmission lines
c. how signals are conveyed from one point to another
d. how transmission lines are arranged
2. As we understand from the text, .......…. .
a. a rhombic antenna is the most popular antenna used in TV systems
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b. a folded-dipole antenna may also be employed where unbalanced
properties are needed
c. waveguides are not normally used below 1 GHz
d. coaxial lines are not normally used between 1 and 18 GHz
3. Paragraphs 2, 3, and 4 mainly describe ……… .
a. the balanced and unbalanced transmission lines
b. the fundamentals of transmission lines
c. practical transmission lines for use in audio-frequency applications
d. practical transmission lines manufactured in different forms
4. It is true that ………. .
a. flashover due to a high-voltage gradient has no effect on high-power
b. flashover may be reduced due to the high reactive property of nitrogen
c. a high-power cable of small cross-sectional dimension can withstand
serious flashover
d. a high-power cable must be made so as not to give up under flashover
5. As we understand from Figure 17-2, ………. .
a. all the quantities shown cause equal problems throughout the length
of the line
b. all the quantities shown are proportional to the length of the line
c. resistance along the line occurs between the two wires in the cable
d. shunt conductance along the line is due t high resistivity of wires in
the cable
C. Answer the following questions orally.
1. What are the two types of transmission lines commonly used?
2. What is the use of parallel-wire line?
3. What are the advantages of rigid cables over the flexible one?
4. What does a rigid air-dielectric line consist of?
5. What is a spacer?
6. What comprises a flexible air-dielectric cable?
7. What causes the capacitance along the line?
8. How are the quantities L, R, C, and G, considered at radio frequencies?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. What is formed by two coaxial conductors is
a. a parallel-wire line
b. a directional-power relay
c. a coaxial line
d. a signal carrier
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2. One type of ………line is the two-wire open line which is sometimes
used as a transmission line between antenna and transmitter or antenna
and receiver.
a. rigid air-dielectric
b. flexible air-dielectric
c. parallel
d. coaxial
3. The electric and magnetic fields in the two-wire parallel line extend
into space for relatively great distances, and …………losses occur.
a. transmission
b. power
c. reflection
d. radiation
4. Any one of a class of antennas producing the radiation pattern
approximating that of an elementary electric dipole is known as ………
a. rhombic
b. grounded
c. dipole
d. quarter-wave
5. The property of a system of conductors and dielectrics that permits the
storage of electrically separated charges when potential differences exist
between the conductors is referred to as………… .
a. resistance
b. capacitance
c. inductance
d. conductance
B. Fill in the blanks with the appropriate form of the words
1. Flexible
a. Concentric cables may be made, with the inner conductor consisting
of ………. wire insulated from the outer conductor by a solid,
continuous insulating material.
b. Early attempts at obtaining ……….. employed the use of rubbed
insulators between the two conductors.
2. Shield
a. The ………. pair consists of parallel conductors separated from each
other and surrounded by a solid dielectric.
b. The conductors are contained within a copper braid tubing that acts
as a …………..
c. The fields are confined to the space between the two conductors;
thus, the coaxial line is a perfectly ……….line.
3. Ground
a. The may be connected to ground without affecting
operation of the device.
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b. No electric or magnetic fields extend outside of the………. conductor
in a coaxial line.
c. A ground bus is a bus to which the ………. from individual pieces of
equipment are connected, and that, in turn, is connected to ground
at one or more points.
d. A ground cable band is used for ………. the armor or sheaths of
cables or both.
4. Space
a. Space charge is the electric charge in a region of ………. , due to the
presence of electrons and/or ions.
b. The direct ………. wave is basically limited to so-called line-of-sight
transmission distances.
c. A ……….shaft is a separate shaft connecting the shaft ends of two
d. A ……….pulse or space is the signal pulse that, in direct-current
neutral operation, corresponds to a circuit open, or no current
5. Break
a. The length of a multiple ……….is the sum of two or more breaks.
b. A motor develops the breakaway torque to………. away its load from
rest to rotation.
c. Breaking capacity is the current that the device is capable of ……. at
a stated recovery voltage under prescribed conditions of use
and behavior.
C. Fill in the blanks with the following words.
The impedance of the cable, termed the ……….or surge impedance Z0, is
considered to…….,independent of the cable length and the……. Frequency.
This consideration is valid when the…….. is properly terminated and when
the operating………. is above a few tens of kilohertz, but…… a few gigahertz
D. Put the following sentences in the right order to form
paragraph. Write the corresponding letters in the boxes
a. A quantitative indication of the nature of a particular standing wave is
given by the standing-wave ratio (SWR).
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b. When voltage and current waves are reflected on a line due to a
discontinuity, standing waves are produced.
c. It is of the nature of this pattern that there are points of maximum and
minimum values.
d. Standing waves are the result of the summing of instantaneous values
of incident and reflected waves at every point along a line.
e. The standing-wave ratio is defined as the ratio of the maximum value
of a wave to its minimum value.
f. The summing process produces a pattern of variation ( the standing
wave) along the line.
Section Two: Further Reading
A source has no way of knowing whether a line is infinite or finite when it
begins to supply current and voltage waves to the line. If the line is
terminated (connected to a load) in a resistance whose value is equal to Z0,
the voltage and current waves will 'enter' that resistance and be dissipated.
The energy that the waves represent will be taken off the line by the
terminating device (the resistance); none of the energy will be returned to the
On the other hand, if the line is simply an open line of finite length,
something must happen to the waves when they reach the end of the line.
Since there is nothing connected to the line to absorb them, they will be
reflected back from the end of the line and will travel along the line toward
the source. On the line there will now be voltage and current waves coming
from the source, and voltage and current waves traveling back from the-end of
the line. The waves from the source are called incident waves, those reflected
from the end are called reflected waves. As with any ac voltage or current, the
two sets of waves will combine phasorally at each point along the line-the
incident voltage wave with the reflected voltage wave, incident current wave
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with reflected current wave. As a result of the waves combining, there will be
established on the line, patterns of voltage and current variations. It happens
that these patterns do not move or travel. These are, therefore, standing waves
and are known by that name.
The standing waves of voltage and current of an open line are depicted
in Figure 17-4. The points in a standing-wave pattern where a voltage or
current is a maximum are called loops, the points where values are minimums
are called nodes.
Distance from open
Is open
Figure 17-4. Standing Waves on an Open Line.
You will observe that on an open line, the voltage standing wave has a
loop at the open end, and the current standing wave has a node at that end.
This is as we would expect from the theory of ordinary circuits-the voltage is
maximum across an open, the current is zero at the open.
Let us imagine that we have all the facilities for exciting a transmission
line with RF energy and for carefully measuring transmitted and reflected
power on the line. The line, a parallel-wire line, is open and is slightly longer
than one wavelength at the exciting frequency. Excited, it develops standing
waves because it is neither infinite in length nor terminated in a load that will
remove all of the RF energy transmitted. If not removed, energy reaching the
open end of the line will be reflected. The reflected energy, measured at the
sending end, will be equal to the energy fed to the line from the source minus
the energy lost during its trip from the generator, along the line to the open
end, and back.
We expect some energy to be lost-the usual 12R loss due to current
flowing in the resistance of the conductors, and a much smaller loss in the
dielectric between the conductors. We can predict these losses quite accurately,
however, using facts about conductor size, measured current, and so on.
We are puzzled, therefore, when we discover that the reflected energy, as
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measured at the sending end of the line, is significantly less than the
transmitted energy minus the predicted losses. What can explain the
greater-than-expected loss of energy? The answer is radiation! A small but
significant amount of RF energy has simply left the transmission line and is
traveling away from it. This phenomenon is the one that makes possible all
wireless communication. An antenna is a device that enhances the process of
radiation of RF energy from a system.
The radiation of electrical energy is very common. The science of
measuring and predicting radiation is highly developed. However, because of
its nature-it is silent, invisible, and odorless-it does not lend itself to an
explanation in simple physical terms. The explanation that follows is a
common one. It does not tell the entire story, but is useful in providing a
working understanding of radiation independent of the very elegant but highly
sophisticated mathematical treatments.
To continue with the example of the transmission line of the preceding
paragraph, let us imagine the conductors of the last A/4 of the line being
spread apart slightly, as in Figure 17-5a. The standing waves of voltage and
current produce an electric field between the conductors and a magnetic field
around each conductor. These are represented by the lines of force of Figure
17-5b. Notice the fringing (bowing out) of the electric field lines at the end of
the line. (Fringing of electric field lines is a common phenomenon at the
boundaries of an electric field between two conductors.) Fringing occurs
because field lines running in the same direction exert a repulsive force on
each other. At a field boundary, this force produces the spreading out and
bowing of the lines.
When the electric field lines are the product of an ac voltage, the lines
must be produced, collapsed, and reproduced in a reverse direction at a rate
equal to the frequency of the voltage, We can imagine the lines being sent out
from, and withdrawn to, the conductors. It is useful to theorize that, when the
frequency exceeds approximately 20,000 cycles (40,000 reversals) per second,
the outermost line of the field simply cannot keep up with the process of
reversal. Being unable to return to its conductor, it closes upon itself, forming
a closed loop. This loop is repulsed by the outermost line of force produced
by the next alternation of the ac voltage. That line of force subsequently fails
to make it back to the conductor and forms another closed loop, etc. The
process repeats itself during each cycle. The closed loops are driven farther
and farther away from the conductors. The result is a continuous wave train of
energy being discharged (radiated) and repulsed from the conductors.
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Figure 17-5. (a) Fringing of Electric Field Lines at the End of Transmission
Line; (b) Cross-Sectional View of Twin-Lead Transmission Line With Electric
and Magnetic Field Lines.
A Basic Antenna: The Half-Wave Dipole
Let us return to the picture of energy being radiated from the slightly spread end
of a parallel-wire transmission line. Refer to Figure 17-5 again. What might
we do to maximize radiation, since that is what is desired in an antenna? The
answer is to bend further the final quarter-wavelength of each conductor until
each is at right angles to the line (sec Figure 17-6). Since
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the length of each conductor
bent is A/4, the overall length of
the perpendicular portion of the
line is now 1/2A. The device
thus produced is called a halfwave dipole. The half-wave dipole
is a simple, basic antenna. It is
Figure 17-6. Half-Wave Dipole Antenna Made by
commonly used as a basis for
comparison for more complex Spreading Apart Conductors of 4 of Transmission
Information about standing
waves on transmission lines can be transferred to the half-wave dipole
antenna. It is like an open line. The distance from the tip of each pole to the
center feed point is  /4. Hence there will be voltage standing-wave loops at
the tips of the dipole (like the loop at the end of
an open line). There will be a voltage wave node
 4 away, at the center feed point. Similarly,
there will be current standing-wave nodes at the
tips, and a current loop  4 away, at the center
feed point (see Figure 17-7). The standing-wave
Figure 17-7. Voltage and Curpattern just
described is exactly the one
rent Standing Waves on Halfneeded
Wave Dipole Antenna.
to maximize radiation.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. The second paragraph mainly describes ............. along an open line
of finite length.
a. the patterns of voltage and current variations
b. the voltage and current waves traveling
c. incident waves, reflected waves, and standing waves
d. incident waves variations at each point
2. The mechanisms of an antenna is based on the process of..……….. .
a. radiation
b. reflection
c. power dissipation
d. line excitation
3. Paragraph 5, 6, and 7 mainly discuss ........... .
a. transmission lines excited by RF energy
b. parallel-wire lines longer than one wavelength
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c. basic properties of conductors
d. basic principles of antennas
4. It is true that ……… .
a. the fringing of electric field lines formed at the end of the line is
compatible with the fringing of magnetic field lines around each
b. the electric field lines produced by an ac voltage lead to a continuous
wave train of energy radiated and repulsed from the conductors
c. the bowing of electric field lines may occur at any point along the line
d. the bowing of electric lines occurs because of the field lines
running in opposite directions
5. We may conclude from the text that ……… .
a. antennas are employed for only the generation of electromagnetic
b. antennas receive electromagnetic waves and convert them into RF
c. an antenna is a passive device; that is, it cannot add any energy to
a signal that has been fed to it for processing
d. an antenna is identical to a circuit containing a transistor; that is,
it adds energy to the signal it is processing
6. It is true that ………. .
a. if we have maximum radiation from an antenna, all energy applied to
it will be converted to electromagnetic energy and radiated
b. If an antenna is made up of a parallel-wire line with two quarterwave sections, the electromagnetic energy radiated will be maximized
c. the distance between the conductors of a parallel-wire line causes
the formation of loops at the tips of the dipole
d. the distance between the conductors of a parallel-wire line causes the
radiated energy to be maximized
B. Write the answers to the following questions.
1. How do you explain the loops and nodes on an open line?
2. What are the values of voltage and current at the end of an open line?
3. What causes energy losses of a transmission line?
4. What kind of energy travels away from an open transmission line?
5. What is wireless communication based on?
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6. What are the characteristics of an open-ended line?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Traveling Voltage and Current Waves
Let us imagine that a parallel-wire transmission line is connected through a
switch to a source of RF voltage, as in Figure 17.8a. Let us imagine,
that the line is infinite in length. At the moment the switch is
closed the effect
of an electrical disturbance begins to be felt on the line.
The effect is that of a
Radio-frequency sinusoidal voltage.This effect travels down(or along) the line
at approximately the speed of light. (The speed, may be somewhat less than
the speed of light depending on the exact nature of the line.) At any given
instant, because of the sinusoidal nature of the voltage, at some points along
the line the voltage will be zero volts, at other points it will be maximum
positive volts. At still other points the voltage will be equal to the peak
negative amplitude, or it will be equal .to anything in between these values. In
other words, at any given instant there is a pattern of sinusoidal voltage
variation along the line. And this pattern is traveling away from the source.
We say that a voltage wave is traveling down the line from the source. The
idea of a voltage wave is shown in Figure 17-8b.
Since the line is imagined to be one of in(infinite length, its input
impedance will be equal to its Z The source will supply a current to the line
with a value given by I  V Z . Because Z is resistive il1 its nature, this
current will be sinusoidal and in phase with the source voltage. The current
effect will travel down the line just at the voltage effect did. We say that there
is a current wave traveling down the line. The current wave is depicted in
Figure 17-8c.
In summary, when a line of infinite length is connected to a source of ac
voltage there is produced on the line a traveling voltage wave and a traveling
current wave. These waves travel away from the source toward the opposite
end of the line. The current wave is in phase with the voltage wave at every
point along the line; its amplitude is determined by I  V Z .
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Of source
To infinity
Figure 17-8. (a) Transmission Line and RF Source; (b) Traveling Voltage Wave
on Transmission Line; (c) Traveling CurrentWave on Transmission Line
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1. circumstance
2. concentric cable
3. conductance
4. convenience
5. corrugate
6. dielectric
7. dimension
8. flashover
9. flexible air antenna
10. folded dipole antenna
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12. gradient.
13. grounded antenna
14. half-wave dipole
15. incident wave
16. leakage
17. node
18. odorless
19. open-ended line
20. parallel-wire line
21. quarter-wave antenna
22. reflected wave
23. rhombic antenna
24. rigid cables
25. sheath
26. standard-wave ratio (SWR)
27. standing wave
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Section One: Reading Comprehension
Any system of wires may be used as a transmission line, but the simplest
arrangements are invariably preferred in practice. Thus parallel-wire and
coaxial lines are by far the most common. In a similar way, a pipe with any
sort of cross section could be used as a waveguide, but the simplest cross
sections are preferred. Accordingly, waveguides with constant rectangular or
circular cross sections are normally employed, although other shapes may be
used from time to time for special purposes. As with regular transmission
lines, so in waveguides, the simplest shapes are the ones easiest to manufacture, and the ones whose properties arc simplest to evaluate. A rectangular
waveguide is shown in Figure 18-1, as is a circular waveguide for comparison.
In a typical setup, there may be an antenna at one end of a waveguide and
some form of load at the other end. The antenna generates electromagnetic
waves, which travel down the waveguide to be eventually received by the load.
It is seen that the waves arc truly guided.
Figure 18-1. Waveguides, (a) Rectangular; (b) Circular.
The walls of the guide are conductors, and therefore reflections from
them take place. It is of the utmost importance to realize that conduction of
energy takes place not through the walls, whose function is only to confine this
energy, but through the dielectric filling the waveguide, which is usually air. In
discussing the behavior and properties of waveguides, it is necessary to speak of
electric and magnetic fields, as in wave propagation, instead of voltages and
currents, as in transmission lines. This is the only possible approach, bu t it
does make the behavior of waveguides more complex to grasp.
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Because the cross-sectional dimensions of a waveguide must be of the
same order as those of a wavelength, use at frequencies below about 1 GHz is
not normally considered, unless special circumstances warrant it.
Both waveguides and transmission lines can pass several signals
simultaneously, but in waveguides it is sufficient for them to be propagated in
different modes to be separated. They do not have to be of different
frequencies. Again, a number of waveguide components are similar if not
identical to their coaxial counterparts. These components include stubs,
quarter-wave transformers, directional couplers and taper sections. Indeed, the
operation of a very large number of waveguide components may best be
understood by first looking at their transmission-line equivalents.
A major problem with twin-lead transmission lines at higher frequencies
is that the amount of direct radiation from such lines increases with the
frequency of the signal being transmitted. The result is that a twin-lead
transmission line radiates virtually all of the energy it is carrying, and
transmits little, if any, to a load at frequencies above several hundred
megahertz. The problem of energy loss due to radiation is almost totally
eliminated with coaxial lines and waveguides because these forms of
transmission lines ‘enclose’ the signal and prevent its radiation.
A second source of energy loss for both parallel-lead and coaxial
transmission lines is in the dielectric that supports the separation of the
conductors. This is called dielectric loss. Although, theoretically, there is no
current flow in an insulator, a dielectric, there is some current flow in actual,
practical dielectrics, and there is dissipation. Of course, this dissipation is
extremely small. However, it increases with frequency. Again, at very high
frequencies, an energy loss becomes consequential. Because waveguides are
completely hollow and in most cases filled with air, dielectric loss is virtually
A third form of energy loss in transmission lines is in the I R heating of
the conductors of the line. Heating or ‘copper’ loss is directly proportional to
the resistance of a conductor, for a given current. And the resistance of
conductors of RF energy increases with frequency! This is the result of the
phenomenon called ‘skin effect’. As the frequency of a current increases, it
‘travels’ more and more on the surface of a conductor. The penetration of the
disturbance of electron movement becomes shallower. This means that a
smaller cross section of a conductor is utilized for current flow. And the
consequence of that, in turn, is an increase in the resistance of the conductor
since resistance is inversely proportional to cross-sectional area. Coaxial lines
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represent some improvement over parallel-wire lines in the matter of heating
loss since the conduction area of the outer conductor is significantly larger
t h a n t h a t of the inner conductor. However, a waveguide has a major advan tage in this regard: the inner (or one) conductor is completely eliminated; and
the conduction area of the inner surface of the guide is significantly larger
than t hat of the coaxial line.
The use of waveguides is not all gravy. In comparison w i t h ot her forms
of transmission lines, they are difficult and expensive to install. The skills
required for installation are more like those of a plumber than of an
electronics technician, or even of an electrician. Waveguides, in most
instances, are rigid devices. Their routing must be carefully planned. Joints or
connection points must be carefully made to avoid discontinuities in the inner,
reflecting surfaces and the consequent creation of standing waves. Other
forms of transmission lines are relatively flexible and can simply he unrolled
and positioned to conform with almost any surface contour.
Waveguides are more expensive to manufacture. They must be precision
made. Inner surfaces must conform to precise dimensions and be free of burrs,
unevenness, etc., which could disturb the reflection patterns of the guided
waves. Several waveguide sections of various shapes used to accommodate a
variety of routing situations are shown in Figure 18-2.
Figure 18-2. Miscellaneous Waveguide Sections.
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
.....… 1. Compared w ith other shapes, rectangular waveguides are the
most common.
.....… 2. Propagation in rectangular waveguides is too difficult to evaluate
compared with other shapes.
.....… 3. In a waveguide, conduction of energy takes place through the air
filling it.
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.....… 4. In order to evaluate the behavior of waveguides, their electric and
magnetic properties must be considered.
.....… 5. Energy loss due to radiation is increased with coaxial lines.
B. Choose a, b, c, or d which best completes each item.
1. Electromagnetic signals propagated through waveguides ........…. .
a. may be of different modes
b. must be of different frequencies
c. may not have frequencies below 1 GHz
d. cannot be evaluated properly
2. it is true that .......... .
a. waveguides stop radiation loss
b. waveguides have the problem of radiation loss
c. twin-lead and coaxial lines decrease energy loss
d. twin-lead line transmits the energy it carries
3. Waveguides are preferred to twin-lead lines because their .......….. loss is
almost nil .
a. radiation
b. dielectric
c. heating
d. all of the above
4. According to the text, dielectric loss in coaxial transmission lines .......... .
a. decreases with frequency
b. increases with frequency
c. is practically high
d. is virtually nil
5. As the frequency of a current along a conductor increases, .......... .
a. the resistance of the conductor decreases since resistance is directly
proportional to frequency
b. the penetration of the disturbance of electron movement becomes
c. the efficiency of the conductor increases, too
d. the conduction area of the conductor increases, too
C. Answer the following questions orally.
1. Why are waveguides not normally used at frequencies below 1 GHz?
2. What similarities are there between transmission lines and waveguides?
3. What are the disadvantages of waveguides over other transmission
4. What are the skills required for waveguide installation similar to?
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5. What difficulty might arise from rigid-waveguide installation?
6. Why are waveguides expensive to manufacture?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. A state of a vibrating system to which corresponds one of the possible
propagation constants is known as …….. .
a. mode voltage
b. mode purity
c. mode
d. frequency
2. At high frequencies, the I R loss is mainly due to …….. .
a. the skin effect
b. the skin depth
c. the self-inductance
d. the self-impedance
3. For any component of a field, the ratio of the instantaneous value of
……… at one point to that of any other point does not vary with time.
a. a waveform
b. a wave envelope
c. a wavefront
d. a standing wave
4. Transmission lines used for transmitting microwaves often take the
form of completely hollow cylindrical or rectangular tubes called …… .
a. parallel lines
b. coaxial lines
c. waveguides
d. waveforms
5. The time rate at which electric energy is transformed into heat in a
dielectric when it is subjected to a changing electric field is referred to
as ……… .
a. dielectric loss
b. dielectric loss angle
c. dielectric loss factor
d. dielectric loss index
B. Fill in the blanks with the appropriate form of the words
1. Reflect
a. When a mismatch occurs, there is an interaction between the
incident and ……….. waves.
b. When a line is terminated with a short circuit, open circuit, or purely
reactive load, no energy can be absorbed by the load so that total
…….. takes place.
c. The voltage ……… coefficient is defined as the ratio of the complex
electric field strength of thereflected wave to that of the incident
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2. Guide
a. The energy of a …….. wave is concentrated within or near boundaries between materials of different properties.
b. The wavelength in a waveguide, measured in the longitudinal direction is known as ……… wavelength.
3. Dissipate
a. Theoretically, there is no current flow in insulators, however, some
current always flows in practical dielectrics and there is ……… .
b. If a line is terminated in a resistance, the voltage and current waves
will enter that resistance and they will be ………. .
4. Contour
a. Flexible transmission lines can be unrolled and positioned to conform
with any surface………. .
b. In a ……… control system, the controlled path can result from the
coordinated simultaneous motion of two or more axes.
5. Enclose
a. Waveguides ………… the signal and prevent its radiation.
b. An ……… relay has both coil and contacts protected from the
surrounding area.
c. A protective housing used to contain equipment and prevent personnel from accidentally contacting live parts is called an …….. .
C. Fill in the blanks with the following words.
When the width, w, of a waveguide is exactly ………… to half of the
free-space wavelength of a/an ……... (i.e., when A=2w), the wave will bounce
…….side to side in the waveguide and will make no progress down the guide.
The angle of incidence or ………. is zero. The frequency that makes this
condition true is …… the cutoff frequency for the guide. The associated….…
is called the cutoff wavelength. Furthermore, the waveguide will…….
attenuate all frequencies lower than the cutoff …….. (or all wavelengths
longer than the cutoff wevelength). A/An……….. is like a high-pass filter. Its
attenuation is …….. that of a resonant wave trap: it does not …… the energy
in the heating of a conductor, it ……… blocks the passage of the energy.
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D. Put the following sentences in the right order to form a
paragraph. Write the corresponding letters in the boxes
a. Waveguides must be excited by a generator in such a way that waves
capable of being propagated will be produced at the point of excitation.
b. Waveguides can be excited by connecting a generator to cither a probe
or loop inserted in the guide, or through a window usually called an iris
or Slot.
c. Unlike methods of exciting twin-lead and coaxial lines, it is not enough
simply to connect a generator to two points on guide.
d. As we have seen, waveguides are literally ‘guides for electromagnetic
waves’ .
Section Two: Further Reading
Theory of Operation
In many respects, waveguides can be dealt with in ways not unlike those used
for other transmission lines. Matters of characteristic impedance, the need for
impedance matching, etc., are not significantly different for waveguides.
However, learning a few new ideas is required if one wishes to gain an
understanding of how a waveguide transmits energy. This understanding is
useful as a basis for a working knowledge of some of the operating
peculiarities and limitations of waveguides. Conventional explanations of
waveguide theory utilize the concepts of electric and magnetic fields
It is useful to think of a waveguide as a special environment for the
propagation of electromagnetic (EM) waves. The ideas involved are not
significantly different from those we examined in connection with the propagation of such waves from antennas. In fact, waveguides are energized or
excited by a probe which acts very much like an antenna. Energy is removed
from a waveguide by an antenna-like probe (see Figure 18-3).
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Figure 18-3.
Input-Output Coupling fur Waveguide Operation.
For electromagnetic waves, the ‘special environment’ which is a
waveguide is like a tunnel. However, the waves are not able to simply streak
straight down the tunnel like a train going through a tunnel. Rather, the
waves are guided in their motion through the tunnel by bouncing from side to
side of the tunnel. Indeed, although the propagation velocity of the waves
along their zigzag journey is the same as that in free space (the speed of
light), the velocity along the axis of the waveguide is less than the speed of
light. This, of course, is (he result of the actual distance traveled being greater
than the length of the tunnel.
You will recall that electromagnetic waves consist of two inseparable
components: an electric field component designated E, and a magnetic field
component designated H. These components are vector quantities since they
have both magnitude and direction in space. The vectors are always at right
angles to each other. Together, the two vectors define a plane. The direction
of travel (of propagation) of an EM wave is always perpendicular to the plane
of the vectors.
When radiation is emitted from a point source-a small, simple
antenna-the electromagnetic energy travels (is propagated) away from that
source in all directions. Since the source of energy varies in amplitude at a
radio-frequency rate, the intensity of the energy being propagated varies at the
same rate. The result is that in the space surrounding the antenna, the energy
intensify varies in a wave-like pattern. The pattern travels out from the
source. The leading edge of this energy disturbance is called, appropriately, a
wavefront. In the first instant of emission, the wavefront is like a small sphere
surrounding the source. With time the wavefront travels away from the
source, expanding the sphere. The action is like that of a spherical balloon
being inflated-the wavefront corresponds to the surface of the balloon. With
each succeeding alternation of the source, a new wave front is generated, and
so on.
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A wavefront is an equiphase surface. That is, the surface represents all
points of equal intensity of the electric and magnetic fields. Between
successive wavefronts, fields vary in a sinusoidal pattern. As the sphere of the
wavefront gets larger, with greater distance from the source, the wavefront
approaches the nature of a flat plane. A wavefront, or simply, a wave, can
never be perfectly flat. However, in examining the theory of wave propagation
in waveguides it is convenient to think in terms of waves with perfectly flat
fronts. Such a wave is called a uniform plane wave.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. As we understand from the text, waveguides ……… .
a. are based on operational theories completely different from those of
other transmission lines
b. are energized by a probe whose mechanism is similar to an antenna
c. cannot propagate electromagnetic waves as efficiently as other
transmission lines
d. cannot be excited by antenna-like probes
2. It is true that ………. .
a. waves travel along a waveguide in a straight line
b. waves bounce back and forth as they move along a waveguide
c. the actual distance traveled by waves along a waveguide is greater
than the length of the waveguide
d. the actual distance traveled by waves along a waveguide is equal to
the length of the waveguide
3. We can infer from the text that ……… .
a. electromagnetic waves can be sent straight down a waveguide
b. electromagnetic waves are not influenced by the waveguide walls
c. the velocity of propagation in a waveguide can be the same as the
speed of light
d. the electric field, the magnetic field, and the direction of propagation
of an EM wave are mutually perpendicular
4. Coordination of variations in amplitude and intensity of the energy
propagated result in ………. .
a. the energy traveling out from the source in a wave-like pattern
b. the energy traveling away from the source in all directions
c. the uniformity of the waves in free space
d. the concentration of the waves in a direct line
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5. The last paragraph mainly describes …….. .
a. variation of fields between successive wavefronts
b. variations of the sphere of a wavefront
c. the characteristics of wavefront
d. the points on the surface of a wavefront
B. Write the answers to the following questions.
1. What is the function of a waveguide?
2. Why is the propagation velocity of the waves along the axis of the
waveguide less than the speed of light?
3. Why are the electric and magnetic field components of electromagnetic
waves vector quantities?
4. What is a wavefront?
5. How does the wavefront vary traveling away from the source?
6. What is a wavefront compared with?
7. How are new wavefronts generated?
Section Three: Translation Activities
A. Translate the following passage into Persian.
Reflection of Waves From a Conducting Plane
As already discussed, an electromagnetic plane wave in space is transverseelectromagnetic, or TEM; the electric field, the magnetic field and the
direction of propagation are mutually perpendicular. If such a wave were sent
straight down a waveguide, it would not, despite appearances, propagate in it.
This is because the electric field (no matter what its direction) would be
short-circuited by the walls, since the walls are assumed to be perfect
conductors, and thus a potential cannot exist across them. What must be
found is some method of propagation which does not require an electric field
to exist near a wall and simultaneously be parallel to it. This is achieved by
sending the wave down the waveguide in a zigzag fashion, bouncing it off the
walls and setting up a field that is maximum at or near the center of the
guide, and zero at the walls. In this case the walls have nothing to shortcircuit, and therefore they do not interfere with the wave pattern set up
between them; thus propagation is not hindered.
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Two major consequences of the zigzag propagation are apparent. The
first is that the velocity of propagation in a waveguide must be less than in
free space, and the second is that waves can no longer be TEM. The second
situation arises because propagation by reflection requires not only a normal
component but also a component in the direction of propagation for either
the electric or the magnetic field, depending on the way in which waves are set
up in the waveguide. This extra component in the direction of propagation
means that waves are no longer transverse-electromagnetic, because there is
now either an electric or a magnetic additional component in the direction of
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1. accommodate
2. boundary
3. circular
4. coaxial line
5. coefficient
6. cutoff frequency
7. cutoff wavelength
8. dielectric loss
9. dissipation
10. equiphase
11. iris
12. mismatch
13. modepurity
14. penetration
15. perpendicular
16. probe
17. rectangular
18. rigid waveguide
19. skineffect
20. slot
21. sphere
22. transmission
23. twin-lead line
24. uniform plane wave
25. wave front
26. waveguide
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Section One: Reading Comprehension
Optical Communication Systems
An optical communications system is primarily a ‘conventional’ telecommun nications system that utilizes light waves as a carrier in one or more
transmission links. Although lightwaves can carry analog signals, optical
systems are invariably also digital communications systems.
An optical communications system, then, is one in which the transmission link is an optical transmission line instead of a terrestrial metallic
conductor transmission line, or microwave link or satellite link, etc. Such a
system has terminal facilities incorporating typical digital communications
functions: analog-to-digital and digital-to-analog converters, multiplexers and
demultiplexers, carrier generators (light sources, in this case) and modulators,
receivers and demodulators, and so on. What is really new and different to be
learned about an optical communications system is the operation of the optical
transmission line.
A special tube made of glass or plastic for the purpose of guiding light is
called an optical fiber. An optical fiber is a waveguide for light waves. The
term 'fiber' is appropriate because this tube or guide is a slender, thread-like
structure. 'Optical' means that it has to do with light. An optical fiber is able
to guide or conduct light along a path that is not a straight line. It can
accomplish this feat with only a minimum of attenuation of the light. Fibers
with attenuation characteristics of the order of 0.2 dB/km (decibels per
kilometer) have been demonstrated in the laboratory. By comparison, the
attenuation of 19-9auge twisted-wire-pair transmission line (in a multi pair
cable) at voice frequencies is about 0.6 dB/km. Because attenuation on
twisted-wire-pair line increases rapidly with frequency, operation is limited to
approximately 1 MHz.
It is not difficult to conceive of a perfectly straight tube functioning as a
light guide; light could simply shoot down the tube. Optical fibers, however,
are seldom perfectly straight. By what principle are they able to conduct light
around curves? The answer is refraction.
Refraction means bending of a wave-like entity such as light. The
direction of a light wave-a light ray-is bent when the ray passes between two
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media in which the velocities of propagation (of light) are different. You have
experienced this phenomenon if you have ever been puzzled when trying to
locate something under water while looking at it from above the surface of
the water. You will recall that the object (e.g., a bar of soap in a bathtub) was
not where you ‘saw’ it to be. The light rays from the object were refracted as
they left the surface of the water. Light travels faster in air than in water.
Optical fibers guide light by refraction. In simplest form, a light
‘conductor’ could consist of a solid glass or plastic rod surrounded by air (see
Figure 19-1). Since the solid and air have different propagation characteristics,
light rays in the rod would be refracted from the sides of the rod and thus be
guided through it.
For various reasons, optical fibers for communications applications have
a form somewhat more complex than the simple ‘light pipe’ of Figure 19-1. As
shown in Figure 19-2, a typical optical fiber consists of three basic elements: a
central core (a solid rod of glass or plastic) surrounded by a protective coating
of a different material called a cladding, which, in turn, is covered with a
protective sheath. Before proceeding further with the specifics of optical
fibers, let us examine briefly the basic 'rules' of refraction and learn some
terminology commonly used in discussions involving optical communication.
Figure 19-1. Concept of Light Travel
Figure 19-2. Three Basic Parts of
Through a Light Pipe.
an Optical Fiber.
Snell's Law
The performance of light rays in refraction at the boundary of two
light-conducting media is predictable from a principle known as Snell’s law.
Before looking at Snell’s law, however, let us learn the meaning of basic terms
associated with refraction, Refer to Figure 19-3. Observe that Figure 19-3
depicts the boundary between two media. Each medium is characterized by a
property called its refraction index, n. The media in the diagram have indexes
of n 1 and n2 ; n 1 is greater than n2 . The index of refraction of a material is
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inversely proportional to the velocity of propagation of light in the material.
For example, air has an index of refraction of approximately 1 (1.0002914); a
typical n for glass is 1.5: light propagates more slowly through glass than
through air. In Figure 19-3, light travels faster in medium 2 than in medium 1.
Figure 19-3. Refraction and Reflection: (a) Angles of Incidence and Refraction
at Boundary Between Media of Different Indexes of Refraction; (b) Ai = Critical
Angle; (c) Ai > Critical Angle.
Examine Figure 19-3 further. Observe that three conditions of a light ray
encountering a media boundary are shown. The conditions relate to different
angles of incidence. The angle of incidence, A i , of an arriving (incident) light
ray is the angle the ray makes with a line perpendicular to the media
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boundary at the point where the ray meets the boundary. The angle of
refraction, A r , is the angle between the perpendicular to the boundary and the
ray as it continues on its way. There is a consistent relationship between A i
and A r Snell’s law:
The ratio of the sine of the angle of incidence in medium 1 (of two media) to the
sine of the angle of refraction in medium 2 is a constant, K, and equal to the ratio
of the index of refraction n2 of the second medium to that, n 1 , of the first:
sin Ai n2
sin Ar n1
Note from Figure 19-3 that when At is relatively small, as in Figure 19-3
(a), the ray is bent but is able to exit medium 1; that is, it is able to cross the
boundary and continue in the general direction of its original path. In Figure
19-3(c), however, where A-t is quite large, the ray is reflected back into
medium 1; it is not able to escape. Figure 19-3(b) illustrates what is called the
critical angle. The critical angle is the incidence angle which produces a
refraction angle of 90°. When the refraction angle is 90°, the ray neither exits
the first medium nor is reflected back into it. Us direction is along the
When optical fibers are used as transmission lines for light in a
communications system, it is important that they be operated so that most of
the light that is introduced to the fiber remains in it until the destination is
reached. That is, the condition of interest is when the angle of incidence is
greater than the critical angle, the condition of Figure 19-3(c).
Figure 19-4. Light Directions at Input to Optical Fiber
In optical-fiber operation, the angle of incidence is set by the
relationship of the light source to the source end of the fiber. Light from a
typical source travels in all directions. The result is that all three of the
conditions illustrated in Figure 19-3 are likely to occur in the excitation of an
optical fiber, as shown in Figure 19-4. Only those rays that enter the fiber
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parallel to its axis or which have incidence angles greater than the critical
angle will be propagated in the fiber.
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
…… 1. Optical systems are used for both analog and digital communications systems.
…… 2. Satellite and optical links are identical.
…… 3. Optical fibers for communications purposes are as simple as light
…… 4. The index of refraction of a material is directly proportional to the
velocity of propagation of light in the material.
…… 5. It may be concluded from the text that when the refraction angle is
90°, the ray is absorbed
…… 6. The best condition in optical-fiber operation is when the angle of
incidence is greater than the critical angle.
B. Choose a, b, c, or d which best completes each item.
1. According to the text, .......... .
a. an optic fiber is quite similar in appearance to a twisted-wire-pair
b. an optic fiber has the same attenuation characteristics as those of a
twisted-wire-pair cable
c. an optical transmission line performs the same functions as a microwave
link but in a different manner
d. an optical transmission line performs function different from those
performed by a microwave link
2. We may infer from the text that .......... .
a. optical fibers are used for light transmissions in a manner virtually
identical to waveguides at microwave frequencies
b. an optic fiber is a piece of very thin, highly pure glass with the same
refractive index as the outside cladding
c. fiber optic system has fully replaced other communications systems
d. fiber optic system is not capable of taking over communication
traffic handled by satellite links
3. In a fiber-optic communications system, modulators must be used to
..… . . .
a. cause the signal to be carried away from the source
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b. cause the light wave to travel down the fiber
c. convert the optical signal back into an electrical signal
d. impress data or an analog signal on the light beam
4. When the angle of incidence is smaller than the critical angle, the ray
………. .
a. is not able to cross the boundary.
b. crosses the boundary and escapes
c. is not able to escape
d. travels along the boundary
5. As we understand from the text, ……… .
a. the speed of light reduces in materials other than the air and this
reduction results in refraction
b. the speed of light reduces in materials other than the air but it has
nothing to do with refraction
c. the angle of refraction increases as the material through which light
passes becomes denser
d. the angle of incidence of an arriving light ray is directly proportional
to its angle of refraction
C. Answer the following questions orally.
1. What does an optical communications system use as a carrier?
2. What are some of the terminal facilities used in an optical system?
3. What principle are optical fibers based on?
4. What does an optical fiber consist of?
5. What are the angles of incidence and refraction?
6. How is the angle of incidence set in an optical-fiber operation?
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. In ……. two or more messages are simultaneously transmitted over the
same transmission path.
a. frequency modulation
b. multiplex operation
c. amplitude modulation
d. fiber-optic operation
2. The electromagnetic waves just below visible light in frequency are the
infrared waves. They are increasingly used in ………. Communication
a. fiber-optic
c. multiplex printing
b. multiplex radio
d. multiple-tuned
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3. The ratio of the phase velocity in free space to that in the medium is
referred to as refraction ……… .
a. error
b. effect
c. index
d. loss
4. A ……… is logically equivalent to a multiposition selector switch.
a. diplexer
b. multiplexer
c. pilot carrier
d. radio detector
5. The ……... of an arriving light ray is the angle the ray makes with a line
perpendicular to the media boundary at the point where the ray makes
the boundary.
a. angle of incidence
b. angle of refraction
c. critical angle
d. none of the above
B. Fill in the blanks with the appropriate form of the words
1. Refract
a. The speed reduction and subsequent ……… are different for each
b. The ………. index, n, is the ratio of the speed of light in free space to
the speed in a given material.
c. When the refraction angle is 90°, the ……….. ray goes along the
d. Electromagnetic waves traveling from a rarer to a denser medium are
……… toward the perpendicular to the boundary.
2. Bend
a. The amount of ……… provided by refraction depends on the
refractive index of the two materials involved.
b. Refraction causes the light to be ………. .
c. The ratio of ………….. amplitude existing before the introduction of
bend-reducing features to that existing afterward is known as bend
reduction factor.
3. Carry
a. Optical systems may be used to ………… digital signals.
b. The process of extracting the signal information from a modulated
………… wave is called demodulation.
c. The current associated with a carrier wave is called the ………
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Figure 19-5.
be glass while the clad is plastic, or both may be made from plastic materials
of different densities. Fiber optics is founded on the theory of reflection that
results at the interface between two materials of different densities. In metallic
waveguide, the energy is reflected along the guide when one-half wavelength
of energy is shorter than the size of the waveguide. In fiber optics, the energy
will reflect down the glass waveguide when the angle of reflection remains
smaller than a critical angle determined by the ratio of the densities of the
core and clad materials.
The cross section of Figure 19-6 illustrates the construction of a cable
having a glass core 50/ m in diameter with an index of refraction of 1.45. The
Figure 19-6. Light Reflection Inside the Guide From N 1 and N 2
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cladding around the core has an outer diameter of 100 m and an index of 1.3.
The clad always has a lower index of refraction than the core material.
Reflections with zero loss will take place at the interface surface of the core
and clad materials, provided the light energy approaches the interface at an
angle that is less than the critical angle  c .
Light energy traveling through a vacuum will move at a velocity of 3 10 8 m/s.
It is considered to have the same velocity in our atmosphere. As light enters
any transparent medium, it slows down slightly depending on the optical
density of the new material. When comparing the velocity of light in free air
to the velocity of light in the given medium, the ratio is a unitless number
called the index of refraction N:
where Vc = velocity of light in air
Vm = velocity of light in the new medium
The index of refraction is a number larger than 1, which means that light in
any transparent material moves slower than it does in air.
The term refraction identifies a directional change to a ray of light, as
well as a velocity change when light crosses between two materials of different
refractive indexes. Refraction is also dependent on the angle of penetration.
This principle can be demonstrated easily. Figure 19-7 illustrates a person
Figure 19-7. Refractions Due to a Change in Index
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looking over the side of a boat into a clear pond. The images we see are
lightrays that are reflected off of the subject and detected by the eye, so we
can use the lines of sight as representing rays of light.
When the person looks straight down into the water, there is a change in
velocity but no change in direction; the seaweed appears directly below the
boat. As the viewer looks toward the shore, he thinks he sees a fish at A. The
angle of penetration is 81, so the light rays refract (bend) as well as reflect,
which means that the fish is actually at B. When the viewer looks at the rock,
the penetration angle 82 is smaller than 81, and the illusion is greater than
when the fish was viewed. The rock appears to be at location A but is really at
location B. As the viewer gazes closer to the shore line, he finds that he can
no longer see into the pond but rather sees the reflection of the tree on the
shore. This is because the angle of entry has become smaller than the critical
angle  c .
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. A simple fiber-optic system ………. .
a. cannot be considered a transmission system
b. is not compatible with other systems
c. consists of a light-emitting source, a cable, and a receiver
d. consists of two lenses and a sensitive receiver
2. According to the text, ……….. .
a. fiber optics and metallic waveguide have the same mechanism
b. fiber optics are based on a theory different from that of metallic
c. the cores of optical fibers have various refraction indexes across their
d. the cores of optical fibers have no relation with the velocity of light
traveling through them
3. We may infer from the text that ………. .
a. in fiber optics, the energy is reflected along the guide when one-half
wavelength of energy is shorter than the size of the waveguide
b. in metallic waveguide, the energy will reflect down the guide when
the angle of reflection remains smaller than the critical angle
c. the information signal is carried by an unmodulated light beam
radiated from a light-emitting source
d. the information signal carried by a fiber-optic transmission line is in
the form of a modulated light beam
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4. Refraction occurs when waves ………. .
a. pass from one density medium to another
b. traveling in straight paths bend around an obstacle
c. travel along the earth’s surface
d. travel through the atmosphere
5. As we understand from the text, ………. .
a. the angle of entry has no effect on the reflection of the objects in the
b. the angle of entry only affects the reflection of the tree
c. the person looking into the water knows exactly where the objects in
the water are
d. the person looking into the water is unaware of the change in velocity
B. Write the answers to the following questions.
1. What is the function of the receiver in a fiber-optic system?
2. Why should the clad have a different optical density from the core
3. What causes the energy to reflect down the glass waveguide?
4. How does the optical density of any material affect the velocity of light?
5. How can an outside cladding of an optic fiber be of the same material
as the core?
Section three: translation activities
A. Translate the following passage into Persian.
As with other forms of transmission lines, energy injected at the input end of a
fiber-optic line diminishes with distance along the line. Light energy in an
optical line is attenuated by four basic loss mechanisms: scattering,
absorption, loss in connections, and loss due to fiber bending.
Scattering and absorption losses are related in that they are caused
primarily by impurities or flaws in the medium of an optical-fiber core. The
ray may be deflected sufficiently to exit the core and be absorbed by the
cladding or sheath of the fiber. Or, it may simply be absorbed by a particle of
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opaque impurity embedded in the core material. In either case, the energy
involved in the particular event is lost to the system.
Connection losses are the result of light rays encountering imperfections
in the boundaries of the transmission media where two media are joined to
permit the passage of light from one to the other. For example, if there is any
surface roughness on the ends of the fibers where they are joined, some light
will be refracted by the surface imperfections and lost to the system.
Bending losses are the result of energy lost when light waves are
required to make an excessively sharp bend. When analyzing the behavior of
light as a wave phenomenon, we must remember that a wave has a ‘width’
perpendicular to the direction of travel of the wave. When the wave bends
around a corner, the outside edge of the wave must travel faster than the
inside edge: If it doesn't, it isn't bending. (As an analogy, when a column of
marchers goes around a corner, persons in the outside positions must step
faster, or persons in the inside positions must mark time, in order for the line
to remain straight during the turn.) If the bend is too sharp, part of the wave
would have to travel faster than the speed of light, which it obviously cannot
do. The result is that some of the light simply exits the fiber and is lost by
absorption in the cladding or sheath.
The total losses of an optical communications system includes the sum of
all the losses produced by the mechanisms described above. Design for
maximum performance requires attention to assure minimization of each type
of loss.
B. Find the Persian equivalents of the following terms and
expressions and write them in the spaces provided.
1. angle of incidence
2. angle of reflection
3. angle of refraction
4. cladding
5. conventional
6. critical angle
7. demultiplexer
8. denser
9. diplexer
10. index of refraction
11. monochrome
12. multiplexer
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13. optical communication system
14. optical density
15. optical fiber
16. reduction factor
17. refraction index
18. Snell’s Law
19. terrestrial
20. twisted-wire-pair cable
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Section One: Reading Comprehension
The Communications Satellite
The ultimate worldwide communications system will feature a satellite as one
of the major components. Although long-range communications took place
before the age of satellites, the systems suffered from conditions that required
a great deal of effort to overcome. Even today, due to remote location or
surrounding terrain, there are isolated communities that are difficult to reach
by point-to-point communication systems. The satellite is really nothing more
than a radio relay station, but it offers the one advantage that is missing in all
other systems; the capability of a direct line-of-sight path to the earth’s surface.
A satellite travels in space in a direction parallel to the surface of a
planet. It has a forward velocity sufficient to create an outward thrust
(centrifugal force) equal to the gravitational pull of the planet it orbits. There
are three common orbital patterns; the polar orbit, the inclined elliptical
orbit, and the equatorial geosynchronous orbit. The following factors apply
equally to all orbits.
1. The plane of the orbit must pass through the center of the object to
be orbited. For instance, a satellite could not orbit the earth around a
latitude of 42°N.
2. The time to complete one orbit depends on the mass of the vehicle
(as compared to the mass of the earth), the vehicle’s velocity
(dependent on the initial thrust supplied by the rocket engines and
the mass of the payload), and the final orbital altitude.
To place a satellite in a position that appears to be stationary over a
selected location on the earth’s surface means that the vehicle must move in
the same direction as the earth rotates. This final requirement eliminates the
polar orbit. An inclined elliptical orbit could be in a direction and at an
altitude and velocity that would appear stationary relative to a given
longitude, but this orbit shifts its north/south latitudinal position.
The only orbit that meets all of these requirements is the equatorial
geosynchronous orbit. It is approximately 22,000 mi or 35,400 km above the
earth’s surface and in a plane that includes the includes the equator.
Satellite communication allows transoceanic links, and wide bandwidths
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Earth sensor
Sun sensor
Solar array panel
Mechanically despun antenna
Omnidirectional antenna
Control electronics for
mechanically despun
Power supply and
low-level traveling
wave tube for
Upper ring
transponder A
Viscous tube
wobble damper
Receive diplexer
Tunnel diode
amplifiers (stacked)
RF oscillator
traveling wave
tube amplifier
Power supply and
low-level traveling
wave tube f o r
transponder B
Power control unit
End cover (kapton)
Lower ring assembly
Propellant / pressurant tank
Figure 20-1.
are utilized to allow the multiplexing of a number of different signals.
Frequencies used are in excess of 1 GHz. At these high frequencies, the
effects of ionospheric refraction and attenuation are negligible. The
frequencies used range from about 1 GHz up to 30 GHz. The signals received
and subsequently retransmitted by the satellite are at different carrier
frequencies. For example, the Intelsat III satellite shown in Figure 20-1
receives signals (the uplink) at from 5.93 to 6.42 GHz, amplifies, translates
down to 3.705 to 4.195 GHz, and then reamplifies via a TWT output stage to
a 7-W level for transmission back to earth (the downlink). The frequency
translation is to prevent interference between the two signals both at the
ground station and satellite. An electronic system performing the reception,
frequency translation, and retransmission is called a transponder. The total
power consumption for satellite operation is about 150 W. The capacity of
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this satellite is for 1200 duplex voice channels or 4 TV broadcasts or any
combination thereof. The Intelsat VII launched in 1992 handles 90,000 voice
The round-trip distance for a satellite relay is typically 90,000 km. The
total transmission time is about 300 ms. Thus, in a transoceanic telephone
conversation, a 600-ms delay occurs before you hear a reply. Because of this,
care is exercised in the routing of international calls to ensure that no more
than a single satellite hop is utilized. Additionally, special circuitry is
incorporated to reduce delayed echo to reasonable levels.
Part I. Comprehension Exercises
A. Put “T” for true and “F” for false statements. Justify your
…….. 1. Satellites have brought about point-to-point communication systems all around the world.
…….. 2. A satellite positioned in the polar orbit does not move in the
same direction as the earth rotates.
…….. 3. A satellite in a satellite communications system is simply a passive
antenna-type reflector.
…….. 4. The total power consumption for a satellite operation is very high.
…….. 5. Intelsat VII is more complicated than Intelsat III.
B. Choose a, b, c, or d which best completes each item.
1. As we understand from the text, a communications satellite ………. .
a. does not necessarily have to be equipped with highly complicated
transmitters and receivers
b. is placed into synchronous orbit, that is, its position remains fixed
with respect to the earth’s rotation
c. can be stationed at any altitude above the earth’s surface
d. must be placed in an orbit compatible with its weight and velocity
2. It can be concluded from the text that ………. .
a. all satellites do not travel around orbits parallel to the surface of the
b. all satellites are equally energized to have the required velocity
c. the plane of an orbit of the latitude of 42°N does not pass through the
center of the earth
d. an orbit of the latitude of 42°N is a good one for a satellite to be
positioned in
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3. At a height of approximately 22,000 miles, ……… .
a. the satellite’s speed is just right to keep it in synchrony with the
rotation of the earth
b. the satellite’s speed must be 22,000 mph to keep it in synchrony with
the rotation of the earth
c. the satellite vehicle naturally moves from west to east
d. the satellite vehicle naturally completes one orbit per hour
4. The time to complete one orbit depends on ………. .
a. the mass of the vehicle
b. the velocity of the vehicle
c. the final orbital altitude
d. all of the above
5. It is true that ……… .
a. the actual transmission links between a communications satellite and
its several stations utilize the medium of narrow-beam microwave
electromagnetic radiation
b. the transmission from an earth station up to a satellite and the
transmission down from a satellite to earth are called downlink and
uplink respectively
c. the satellite vehicle is equipped with a transponder whose function is
only to receive signals from a transmitting earth station
d. the satellite vehicle is not required to provide a source of energy to
operate the communications equipment that it carries
C. Answer the following questions orally.
1. What advantage does the satellite have over other communication
2. What are the three common orbital patterns?
3. What are the advantages of a geosynchronous orbit over other orbits?
4. What are the characteristics of an inclined elliptical orbit?
5. How does ionosphere affect satellite communication?
6. How is the interference between two signals at the ground station and
the satellite prevented?
7. What has been done to reduce delayed echo in transoceanic telephone
Part II. Language Practice
A. Choose a, b, c, or d which best completes each item.
1. A is a transmitter-receiver facility that transmits signals automatically when the proper interrogation is received.
a. relay
b. receiver
c. transponder
d. transmitter
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must provide a very large number of equivalent individual communications channels.
b. The development of digital communications techniques has made the
utilization of TDM possible.
c. Early satcom systems used only frequency-division multiplexing (FDM)
to achieve more intensive utilization of these expensive facilities.
d. Satellite communications systems are extremely expensive as total
e. New systems are using TDM to increase the information -carrying
capacity of equipment.
f. Ground terminal facilities are also sophisticated and costly.
g. The launch costs-the cost of the launching rocket, fuel, launching
facilities, highly skilled personnel, etc. -are a significant part of the total
cost of a satellite system.
Section Two: Further Reading
Time-Division-Multiplexed Earth Terminal
A block diagram showing only the most basic of details of an earth terminal
for a TDM digital satellite communications system is shown in Figure 20-2.
The diagram is for a system utilizing C-band links: 6-GHz uplink and
20-2GHz downlink. You will observe that the 'sending' side of the terminal
includes a multiplexer for selecting, in turn, each of the incoming information
signals. For purposes of this diagram it is assumed that all incoming
information signals have already been converted to some form of PCM (i.e.,
to digital form). These signals would typically arrive over twisted-wire pair or
coaxial transmission lines from various subscribers located in the vicinity of
the earth terminal.
Each incoming signal is allocated a time slot on the uplink carrier. This
allocation is the result of the action of the multiplexer. An RF carrier is
modulated by the digital signals. After modulation, the carrier frequency is
shifted (converted) to that of the uplink-6 GHz-by means of a heterodynetype frequency converter. It is referred to as the upconverter. The output
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of the upconverter is passed through a bandpass filter (BPF) to ensure that its
bandwidth is properly limited. The signal is then amplified to increase its
energy level to that sufficient for transmission over the radio link between the
earth terminal antenna and the satellite antenna. Amplification is typically by
means of a special electronic device for microwave frequencies. The device is
called a traveling-wave tube, abbreviated TWT.
The antenna for the microwave frequency of the 6-GHz uplink signal is
of the parabolic reflector (or 'dish') type. This antenna confines the radiation
to a relatively narrow beam. A narrow-beam radiation pattern has at least two
significant advantages:
1. It concentrates the radiation so as to increase the ERP (effective
radiated power) in the desired direction. This effect is especially
important for the downlink since the amount of power available to
operate a transmitter in the satellite vehicle is extremely limited.
2. The narrow beam reduces the potential for signals intended for one
satellite from interfering with other, nearby satellites.
Refer again to the block diagram of Figure 20-2. Study the portion of
Figure 20-2.
Elements of Earth Terminal for Satcom System: (a) Uplink
Function; (b) Downlink Function.
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the diagram that refers to the ‘receive’ function of the terminal. You will
observe that from the antenna the signal first passes through a BPF and then
a low-noise amplifier (LNA). These two items are typically incorporated as
part of the antenna assembly. It is important to the success of the system that the
downlink signal receive low-noise amplification immediately off the
antenna. This measure helps to ensure that the amplitude of what may be a
relatively weak signal is boosted before the noise content has become
excessive. From the LNA the signal is transmitted, typically through a section
of waveguide, to the downlink receiver. The downlink receiver is a form of
superheterodyne receiver. The 4-GHz signal is converted in the down convener
to a lower intermediate frequency (IF), amplified, and finally, demodulated.
Remember, the carrier is transporting several information signals by means of
time-division multiplexing. These information signals are now separated in a
demultiplexer and sent on their way over terrestrial (land-based) facilities to
their ultimate destinations.
The satellite transponder portion of a satcom system is represented in
Figure 20-3. In brief, the transponder has a receiver section for receiving the
modulated uplink signal (a 6-GHz signal in the example being illustrated
here). The receiver converts this signal to the downlink frequency (4 GHz in
this example) in the function labeled ‘down converter’ in the diagram. The
signal is filtered and receives some preamplification in a TWT device in the
receiver portion of the transponder. It is then fed to the transmitter section of
the transponder. In the transmitter section it receives further filtering; its
energy level is amplified to the desired level for retransmission. Finally, the
downlink signal is fed to a dish-type antenna which is aimed at a particular
area of the earth’s surface.
Comprehension Exercises
A. Choose a, b, c, or d which best completes each item.
1. It is concluded from the first paragraph that the text is a description
providing only the basic details of a hypothetical satellite .......... .
a. communications system
b. terminal elements
c. transponder
d. waveguides
2. It is true that .......... .
a. the uplink function begins with a narrow beam radiated from a dish
b. the multiplexer function is to separate the information signals
c. the receiving and the sending sides of the terminal consist of
identical elements
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